Physics and Technology For Engineers... 2023

Physics and Technology for Engineers
R. Prasad
Physics and Technology

for Engineers
Understanding Materials and Sustainability
R. Prasad
Aligarh Muslim University
Aligarh, U.P, India
ISBN 978-3-031-32083-5 ISBN 978-3-031-32084-2 (eBook)

https://doi.org/10.1007/978-3-031-32084-2
© The Editor(s) (if applicable) and The Author(s), under exclusive license to Springer Nature
Switzerland AG 2023
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does not imply, even in the absence of a specific statement, that such names are exempt from the relevant
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The registered company address is: Gewerbestrasse 11, 6330 Cham, Switzerland
Dedicated to my wife
Sushma Mathur (12 Nov, 1950–24 May,
2023)
(Beautiful lady fought bravely with cancer

for three years)
Preface
The idea of writing this manuscript originated from my interaction with B.Tech.
students who often complaint of not having a book that covers both the fundamentals
of physics and modern technologies. The need of such a book was also felt by the
group of eminent faculty with whom I was involved in setting papers for competitive
examinations of various technical boards/institutes. It was realised that the available
books/materials on these topics are incomplete, lopsided or too detailed in some
aspects but lacking in others. Moreover, in most of the available books, modern topics
like sustainability and sustainable energy sources are not even touched. With the view
of providing a balanced description of physics of engineering materials and modern
technologies and with the aim of making readers aware of their moral responsibility
towards sustainable development, the present text is developed in textbook format.
The book contains around 220 illustrative figures and some 35 tables.
The book is divided into ten chapters. Chapter 1 starts with the classification
of engineering materials and their important properties. In order to link specific
material properties and their dependence on atomic structure of constituent atoms/
molecules, details of atomic structure and of atomic/ionic/molecular bonding are
provided in this chapter. Chapter 2 discusses electrical behaviour of solids including
superconductors and associated physics. Chapter 3 details the origin and behaviour
of magnetic materials and types of magnetism along with the fabrication of materials
with desired magnetic properties. Important topics of modern physics, like discovery
and properties of X-rays, dual nature of matter and instances of the failure of classical
physics based on Newtonian Mechanics and Maxwell’s theory of Electromagnetic
radiations, and their successful resolution in the frame work of quantum approach
are discussed in Chap. 4. Basics of quantum mechanics, particularly of Schrodinger
approach is provided in Chap. 5. Some simple one-dimensional problems of potential
wells and of barrier transmission are discussed in this chapter. Since most micro-
and macrosystems consist assemblies of large number of identical particles, their
behaviour is generally predicted by the laws of statistics. Since microsystems obey
quantum mechanics and have discrete energy levels, the appropriate statistics that
may be applied to these systems is quantum statistics. Laws of quantum statistics,
macro- and microstates, a prior equal probability of all microstates associated with a
vii
viii Preface
given macrostate, etc. are discussed in Chap. 6. Technique of optical fiber communi-
cation is discussed in Chap. 7, while details of laser technology and its applications
are detailed in Chap. 8. Chapter 9 gives detailed description of nanomaterials, their
advantages, reasons behind their special properties, fabrication techniques for nanos-
tructures, membranes, sheets, tubes, etc. Chapter 10 is special as it discusses details
of the concept of sustainability, as applied to different fields like the social, economic,
environmental and sustainable energy sources. Methods conducive to a sustainable
development that may be adopted by individuals, by socio-economic groups, cluster
of groups, etc. are discussed in this chapter. It is expected that after going through
this chapter, a reader will become aware of his/her social responsibility towards
sustainable development and engineers in particular will participate in generating
sustainable energy sources so that the benefits of natural resources are left largely
undiminished for the future generations.
Chapters of the book have the following special features:
1. The objective of the chapter is spelled out at the very beginning.
2. Sufficient number of self-assessment questions (SAQ), probing the understanding
of the reader, is uniformly distributed over the text of each chapter. A serious
reader is expected to satisfy himself/herself by answering these questions before
proceeding further.
3. Solved examples are included, wherever required, to illustrate the technique of
problem solving.
4. Problems with answers are provided at the end of each chapter.
5. Large number of short answer questions are included at the end of each chapter.
6. Since most of competitive examinations are based on multiple choice questions,
sufficient number of multiple choice questions (MCQ) with answers is provided
at the end of each chapter. A special feature of these MCQ is that in some cases
more than one alternative may be correct, and therefore all correct alternatives
must be marked for complete answer of the question.
7. Some long answer questions are also provided at the end of each chapter.
8. Each topic of the text is started from the very basics and is developed to the
desired level; therefore, no other book or material is required for reading this
text.
It is expected that the book will prove useful for readers.
I shall very much appreciate receiving feedback from readers on the following
e-mail address:
Rpm166@rediffmail.com
Kellyville, NSW, Australia R. Prasad

Acknowledgements
Present book, like my earlier publications, is the result of encouragement and support
extended by my students, members of my research group and colleagues. I would
like to express my gratitude to all my students, members of my research group and, in
particular, Prof. B. P. Singh (Former Chairman, Department of Physics) who showed
great enthusiasm for the book project. While he provided some initial material,
his involvement in developing the text was limited due to other academic commit-
ments. Nevertheless, I would like to record my sincere appreciation to Prof. Singh’s
unwavering support, encouragement and the valuable discussions we frequently had.
Support from Prof. Manoj K. Sharma, Prof. Sunita Gupta, Prof. M. M. Musthafa (Ex-
chairman Department of Physics, Calicut University), Dr. Pushpendra P. Singh, Dr.
D. P. Singh, Dr. Abhishek Yadav, Dr. Unnati, Dr. Mohd. Shuaib and Dr. M. Shariq
Asnain is thankfully acknowledged.
I wish to thank the Aligarh Muslim University, Aligarh, India, and my colleagues
at the Physics Department and at the Department of Applied Physics, Z. H. College
of Engineering and Technology, AMU, Aligarh, with whom I passed more than forty
years of my active life.
Last but not the least, I wish to thank all the members of my family, in particular
my wife Sushma who in spite of being seriously ill, extended all possible support
and encouragement for the completion of the project.
Kellyville, Australia R. Prasad
ix
Contents
1 Engineering Materials, Atomic Structure and Bounding . . . . . . . . . . 1

1.1 Classification of Condensed Matter . . . . . . . . . . . . . . . . . . . . . . . . 1
1.1.1 Metals . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2
1.1.2 Ceramics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 17
1.1.3 Polymers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
1.1.4 Composites . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 24
1.2 Atomic Structure . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 29
1.2.1 Elements of Atomic Structure . . . . . . . . . . . . . . . . . . . . 29
1.2.2 Arrangement of Electrons in Atom . . . . . . . . . . . . . . . . 30
1.2.3 Shape and Orientation of Orbitals . . . . . . . . . . . . . . . . . 34
1.2.4 Electron Energy Level Diagram . . . . . . . . . . . . . . . . . . . 36
1.2.5 Electron Configuration of Elements . . . . . . . . . . . . . . . 37
1.2.6 Aufbau or Building-Up Principle . . . . . . . . . . . . . . . . . 37
1.2.7 Representing Electron Configuration . . . . . . . . . . . . . . 39
1.2.8 Valence Shell . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 40
1.2.9 Some Anomalous Electron Configurations . . . . . . . . . 42
1.3 Bonds Between Atoms and Ions . . . . . . . . . . . . . . . . . . . . . . . . . . . 43
1.3.1 Electronegativity . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 43
1.3.2 The Octet Rule . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 46
1.3.3 Classification of Bonding . . . . . . . . . . . . . . . . . . . . . . . . 46
2 Electrical Behaviour of Condensed Matter . . . . . . . . . . . . . . . . . . . . . . . 61
2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 61
2.2 Electron Energy Band Theory . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 64
2.3 Insulators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 66
2.4 Semiconductors . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 68
2.4.1 Intrinsic Semiconductors . . . . . . . . . . . . . . . . . . . . . . . . 68
2.4.2 Covalent Band Picture of Intrinsic
Semiconductor . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 76
2.4.3 Doped or Extrinsic Semiconductors . . . . . . . . . . . . . . . 79
2.4.4 Doping Technology . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 81
xi
xii Contents
2.4.5 n and p Type Semiconductors . . . . . . . . . . . . . . . . . . . . 85

2.4.6 Compensated Semiconductor . . . . . . . . . . . . . . . . . . . . . 90
2.4.7 Degenerate and Non-degenerate
Semiconductors . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 90
2.4.8 Direct and Indirect Semiconductor . . . . . . . . . . . . . . . . 90
2.4.9 Compound Semiconductors . . . . . . . . . . . . . . . . . . . . . . 91
2.4.10 Current Flow in Semiconductor . . . . . . . . . . . . . . . . . . . 92
2.4.11 Temperature Dependence of Semiconductor
Resistivity . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 95
2.4.12 Theoretical Calculation of Carrier
Concentration in a Semiconductor . . . . . . . . . . . . . . . . 96
2.4.13 Hall Effect . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 99
2.4.14 p–n Junction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 101
2.4.15 Some Formulations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 108
2.5 Conductors . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 113
2.5.1 Semimetals and Half Metals . . . . . . . . . . . . . . . . . . . . . . 114
2.6 Superconductor . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 116
2.6.1 Background . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 117
2.6.2 BCS Theory of Superconductivity . . . . . . . . . . . . . . . . . 128
3 Magnetic Materials . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
3.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 137
3.2 Electric Current and Magnetic Field . . . . . . . . . . . . . . . . . . . . . . . 139
3.3 Magnetic Dipole Moment . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 141
3.4 Magnetic Moment of a Charged Particle Moving
in a Circular Orbit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 143
3.4.1 Classical to Quantum Mechanics . . . . . . . . . . . . . . . . . . 145
3.5 Magnetic (Dipole) Moment of Electron . . . . . . . . . . . . . . . . . . . . . 146
3.6 Magnetic Behaviour of Solids . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 150
3.6.1 Magnetic Induction B and Magnetic Field H . . . . . . . 151
3.7 Classification of Magnetic Materials . . . . . . . . . . . . . . . . . . . . . . . 154
3.7.1 Diamagnetic Materials . . . . . . . . . . . . . . . . . . . . . . . . . . 154
3.7.2 Paramagnetic Materials . . . . . . . . . . . . . . . . . . . . . . . . . . 158
3.7.3 Ferromagnetic Materials . . . . . . . . . . . . . . . . . . . . . . . . . 164
3.7.4 Antiferromagnetic and Ferrimagnetic Materials . . . . . 176
3.8 Permanent Magnetic Materials . . . . . . . . . . . . . . . . . . . . . . . . . . . . 182
4 X-rays, Dual Nature of Matter, Failure of Classical Physics
and Success of Quantum Approach . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 191
4.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 191
4.2 Discovery, Production and Properties of X-rays . . . . . . . . . . . . . . 192
4.2.1 Production of X-rays . . . . . . . . . . . . . . . . . . . . . . . . . . . . 192
4.2.2 Continuous X-rays . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 194
4.2.3 Characteristic X-rays . . . . . . . . . . . . . . . . . . . . . . . . . . . . 198
4.2.4 Mosley’s Law . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 201
4.2.5 X-ray Diffraction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 205
Contents xiii
4.2.6 Some Application of X-rays . . . . . . . . . . . . . . . . . . . . . . 208

4.3 Dual Nature of Matter . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 210
4.3.1 Davisson and Germer Experiment . . . . . . . . . . . . . . . . . 212
4.4 Some Examples of the Failures of Classical Approach
and Success of Quantum Approach . . . . . . . . . . . . . . . . . . . . . . . . 222
4.4.1 Stability of the Atom and the Nature of Atomic
Spectra . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 222
4.4.2 Photoelectric Effect . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 224
4.4.3 Quantum Theory of Photoelectric Effect . . . . . . . . . . . 231
4.4.4 Work Function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 234
4.4.5 Residual Atom after the Emission
of Photoelectron . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 234
4.5 Blackbody Radiations and Their Energy Distribution . . . . . . . . . 235
4.5.1 Wien’s Displacement Law . . . . . . . . . . . . . . . . . . . . . . . 237
4.5.2 Failure of Wien’s Distribution Law . . . . . . . . . . . . . . . . 238
4.5.3 Rayleigh and Jean’s Distribution Law . . . . . . . . . . . . . 239
4.5.4 Failure of Rayleigh–Jeans Distribution . . . . . . . . . . . . . 240
4.6 Quantum Theory of Blackbody Radiations . . . . . . . . . . . . . . . . . . 240
4.7 Compton Scattering of Gamma Rays . . . . . . . . . . . . . . . . . . . . . . . 242
4.7.1 Compton Wavelength . . . . . . . . . . . . . . . . . . . . . . . . . . . 245
4.7.2 Compton Scattering by the Whole Atom . . . . . . . . . . . 245
4.7.3 Photon Interactions with Matter . . . . . . . . . . . . . . . . . . . 246
4.7.4 Some Applications of Compton Scattering . . . . . . . . . 247
4.8 Specific Heat of Solids . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 247
4.8.1 Dulong–Petit Law . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 247
4.8.2 Obtaining Dulong–Petit Law on the Basis
of Classical Physics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 248
4.8.3 Problems with Dulong–Petit Law . . . . . . . . . . . . . . . . . 249
4.9 Quantum Approach to Atomic Specific Heat of Solids . . . . . . . . 249
4.9.1 Einstein’s Theory for Specific Heat of Solids . . . . . . . 250
4.9.2 Investigating the Temperature Dependence
of Einstein’s Equation . . . . . . . . . . . . . . . . . . . . . . . . . . . 251
4.9.3 Drawbacks of Einstein’s Model . . . . . . . . . . . . . . . . . . . 252
4.9.4 Debye Theory of Atomic Specific Heat . . . . . . . . . . . . 253
4.9.5 Debye Temperature θD . . . . . . . . . . . . . . . . . . . . . . . . . . 257
5 Introduction to Quantum Mechanics . . . . . . . . . . . . . . . . . . . . . . . . . . . . 267
5.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 267
5.2 Postulates of Quantum Mechanics . . . . . . . . . . . . . . . . . . . . . . . . . 268
5.2.1 What Does Wavefunction Represent? . . . . . . . . . . . . . . 269
5.2.2 Properties of the Acceptable Wavefunction . . . . . . . . . 270
5.3 Observables and Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 270
5.4 Time Evolution of a Quantum Mechanical System . . . . . . . . . . . 271
5.4.1 Schrodinger Time-Dependent Equation . . . . . . . . . . . . 271
5.4.2 Some Properties of Schrodinger Equation . . . . . . . . . . 272
xiv Contents
5.5 Time-Independent Schrodinger Equation . . . . . . . . . . . . . . . . . . . 273

5.6 About Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 274
5.6.1 Null Operator (O) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 275
5.6.2 Unity or Identity Operator ( Iˆ) . . . . . . . . . . . . . . . . . . . . 275
5.6.3 Linear Operator . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 275
5.6.4 Hermitian Conjugate and Hermitian Operator . . . . . . . 275
5.6.5 Anti-hermitian Operator . . . . . . . . . . . . . . . . . . . . . . . . . 276
5.6.6 Inverse Operator ( Â−1 ) . . . . . . . . . . . . . . . . . . . . . . . . . . 276
5.6.7 Unitary Operator (Û ) . . . . . . . . . . . . . . . . . . . . . . . . . . . . 276
5.6.8 Some Properties of Hermitian Operators . . . . . . . . . . . 276
5.6.9 Algebra of Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . 277
5.6.10 Operators for Some Dynamical Variables . . . . . . . . . . 279
5.7 Measurement of a Dynamical Variable in Quantum
Mechanics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 283
5.7.1 Expectation Value of a Dynamic Variable . . . . . . . . . . 286
5.8 Some One-Dimensional Problems . . . . . . . . . . . . . . . . . . . . . . . . . 289
5.8.1 Energy States: Bound and Scattering States . . . . . . . . . 289
5.8.2 Quantum Mechanical Description of a Free
Particle . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 291
5.8.3 Particle in a One-Dimensional Asymmetric
Infinite Potential Well . . . . . . . . . . . . . . . . . . . . . . . . . . . 294
5.8.4 Potential Barrier and Tunnelling . . . . . . . . . . . . . . . . . . 300
5.9 Heisenberg Uncertainty Principle . . . . . . . . . . . . . . . . . . . . . . . . . . 308
5.10 Correspondence Principle and Ehrenfest’s Theorem . . . . . . . . . . 309
6 Quantum Statistics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 317
6.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 317
6.2 Application of Quantum Statistics (Statistical Mechanics)
to an Assembly of Non-interacting Particles . . . . . . . . . . . . . . . . . 318
6.3 Energy Levels, Energy States, Degeneracy and Occupation
Number . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 319
6.3.1 Distinguishable and Indistinguishable Particles . . . . . 322
6.3.2 Macrostate . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 323
6.3.3 Microstates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 324
6.3.4 Time Evolution of an Assembly . . . . . . . . . . . . . . . . . . 325
6.3.5 Postulate of Equal a Prior Probability of All
Microstates . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 327
6.4 Quantum Statistical Probability of a Macrostate . . . . . . . . . . . . . 327
6.4.1 System Properties and Average Occupation
Number . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 328
6.5 The Bose–Einstein Statistical Distribution . . . . . . . . . . . . . . . . . . 328
6.6 The Fermi–Dirac Statistical Distribution . . . . . . . . . . . . . . . . . . . . 332
6.7 The Maxwell–Boltzmann Statistical Distribution . . . . . . . . . . . . . 334
6.8 Relation Between Entropy and Thermodynamic
Probability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 336
Contents xv
6.9 The Distribution Function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 338

7 Optical Fiber Communication . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 347
7.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 347
7.2 Advantages of Optical Fiber Communication . . . . . . . . . . . . . . . . 348
7.3 Basics of Optical Fiber Communication . . . . . . . . . . . . . . . . . . . . 349
7.3.1 Optical Fiber Materials . . . . . . . . . . . . . . . . . . . . . . . . . . 350
7.3.2 Frequently Used Wavelengths in Optical
Transmission . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 351
7.3.3 Principle of Total Internal Reflection . . . . . . . . . . . . . . 351
7.3.4 Types of Fibers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 353
7.3.5 Rays Guided Through Fiber . . . . . . . . . . . . . . . . . . . . . . 356
7.3.6 Meridional and Skewed Rays . . . . . . . . . . . . . . . . . . . . . 356
7.3.7 Acceptance Angle . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 357
7.3.8 Numerical Aperture (NA) . . . . . . . . . . . . . . . . . . . . . . . . 360
7.3.9 The V Parameter . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 362
7.3.10 Attenuation and Dispersion of Optical Signal . . . . . . . 363
7.4 Components of Optical Fiber Network Link . . . . . . . . . . . . . . . . . 369
7.5 Applications of Optical Fiber Transmission . . . . . . . . . . . . . . . . . 373
8 Laser Technology and Its Applications . . . . . . . . . . . . . . . . . . . . . . . . . . 377
8.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 377
8.2 Electromagnetic Radiations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 377
8.3 Interaction of Electromagnetic Radiation with Matter . . . . . . . . . 379
8.4 Einstein Prediction of Stimulated Emission . . . . . . . . . . . . . . . . . 383
8.5 Stimulated (or Induced) Emission of Photons . . . . . . . . . . . . . . . . 386
8.5.1 Population Inversion . . . . . . . . . . . . . . . . . . . . . . . . . . . . 388
8.5.2 Essential Requirements for Laser Action . . . . . . . . . . . 389
8.5.3 Pumping . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 390
8.5.4 Three and Four Level Lasing Schemes . . . . . . . . . . . . . 391
8.5.5 Optical Resonator or Laser Cavity . . . . . . . . . . . . . . . . 393
8.6 Special Characteristics of Laser Light . . . . . . . . . . . . . . . . . . . . . . 399
8.7 Classification of Laser Sources . . . . . . . . . . . . . . . . . . . . . . . . . . . . 405
8.7.1 Solid State Lasers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 406
8.7.2 Dye (Liquid) Laser Source . . . . . . . . . . . . . . . . . . . . . . . 415
8.7.3 Gas Laser Sources . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 417
8.7.4 Excimer Laser . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 422
8.7.5 Mode Locking . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 422
8.7.6 Q-Switching . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 423
8.8 Some Applications of Lasers . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 425
9 Nanomaterials . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 435
9.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 435
9.2 Special Features of Nanomaterials . . . . . . . . . . . . . . . . . . . . . . . . . 437
9.3 Technology Used for the Study of Nanostructures . . . . . . . . . . . . 446
9.4 Techniques of Producing Nanostructures . . . . . . . . . . . . . . . . . . . 452
xvi Contents
9.4.1 Bottom-Up Techniques . . . . . . . . . . . . . . . . . . . . . . . . . . 453

9.4.2 Top-Down Techniques of Fabricating
Nanostructures . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 456
9.4.3 Carbon Nanotubes . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 462
10 Sustainability and Sustainable Energy Options . . . . . . . . . . . . . . . . . . . 473
10.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 473
10.2 Social Sustainability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 474
10.3 Economical Sustainability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 475
10.4 Environmental Sustainability . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 476
10.4.1 Atmosphere . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 477
10.4.2 Mechanism of Trapping Heat by Greenhouse
Gases . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 480
10.4.3 Global Greenhouse Gas Emission by Human
Activities . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 481
10.5 Global Warming . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 483
10.5.1 The Carbon Footprint . . . . . . . . . . . . . . . . . . . . . . . . . . . 485
10.5.2 Reducing and Offsetting Carbon Footprints . . . . . . . . . 485
10.6 Projections on Average Temperature Rise of 1.5 °C
Above Pre-industrial Levels . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 486
10.7 United Nation’s Efforts . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 487
10.7.1 Outlook Scenarios: Computer Model-Based
Scenarios Prepared by IEA . . . . . . . . . . . . . . . . . . . . . . . 488
10.8 Sustainability of Land Mass . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 489
10.9 Sustainability of Water Bodies . . . . . . . . . . . . . . . . . . . . . . . . . . . . 490
10.9.1 Sustainability of River and Other Water Systems . . . . 491
10.10 Some Efforts for Improving the Sustainability
of Environment . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 492
10.10.1 A Unique Fight Against Climate Change;
the Ice Stupa or Artificial Glacier . . . . . . . . . . . . . . . . . 494
10.11 Sustainable Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 495
10.11.1 Units of Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 495
10.11.2 Primary Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 495
10.11.3 Global Energy Production, an Overview . . . . . . . . . . . 496
10.11.4 Electricity: The Most Convenient Form
of Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 497
10.11.5 Cost of Electricity by Source: Cost Metrics . . . . . . . . . 499
10.11.6 Energy Densities Associated with Prevalent
Energy Sources . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 499
10.11.7 Problem with Present Energy Mix . . . . . . . . . . . . . . . . . 501
10.12 Some Clean and Sustainable Sources . . . . . . . . . . . . . . . . . . . . . . . 501
10.12.1 Hydrogen as an Alternative Source of Energy . . . . . . . 502
10.13 Hydrogen Fuel Cell . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 503
10.14 Nuclear Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 506
10.14.1 Drawbacks of Fission Reactor . . . . . . . . . . . . . . . . . . . . 510
Contents xvii
10.14.2 Plus Points of Fission Reactor . . . . . . . . . . . . . . . . . . . . 511

10.14.3 Accelerator-Driven Energy Amplifier . . . . . . . . . . . . . . 512
10.15 Terrain Dependent Renewable Energy Sources . . . . . . . . . . . . . . 514
10.15.1 Geothermal Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 515
10.15.2 Hydroelectric Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . 516
10.16 Wind Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 518
10.17 Solar Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 519
10.17.1 Solar Thermal . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 519
10.17.2 Solar Photovoltaic (PV) Technology . . . . . . . . . . . . . . . 521
10.18 Energy from Ocean . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 523
10.18.1 Tidal Energy . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 523
10.18.2 Ocean Thermal Energy . . . . . . . . . . . . . . . . . . . . . . . . . . 524
10.19 Portable Sources of Sustainable Energy . . . . . . . . . . . . . . . . . . . . 525
10.19.1 Lithium-Ion Battery . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 526
10.19.2 Super Capacitor . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 528
Index . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 533
Chapter 1
Engineering Materials, Atomic Structure
and Bounding
Objective
Classification of condensed matter and its correlation with atomic structure and
chemical bonding are important from the view point of engineers. These topics are
discussed in this chapter in sufficient details. It is expected that after going through
the chapter, the reader will be able to identify the special properties of different
engineering materials and will also be able to correlate these special characteristics
of different materials with their atomic structure, electron configuration and atomic/
molecular bonding. This will go a long way in selecting a proper material for specific
engineering requirements as well as in fabricating materials with desired properties.
1.1 Classification of Condensed Matter
It is known that matter, on the basis of their physical state, may be classified as solids,
liquids, gases and plasma. The first three states are quite well known; however, the
fourth state, plasma, is rather peculiar. At very high temperature atoms of the matter
get ionised forming plasma that contains ionised atoms and electrons in a state of
rapid motion. The flame of a burning candle is a typical example of plasma.
In material science solids are defined as the matter having the property of crys-
tallinity. Crystallinity means molecules, atoms or ions of the matter spaced at regular,
repeating distances and angles from each other in three dimensions. In condensed
matter atoms or molecules or ions almost touch each other, like that in solids; liquids
also show properties of condense matter; however, most of the time liquids do not
have crystallinity: they are amorphous. Some liquids called liquid crystal, as excep-
tion, do exhibit some regularity of structure over comparatively large distances, but
they do not possess this regularity in three dimensions. In supercritical states of
matter, at very high temperature and pressure, matter is essentially in gaseous state
but relative separation between constituent atoms, etc. is of the order of that in solids
© The Author(s), under exclusive license to Springer Nature Switzerland AG 2023 1

R. Prasad, Physics and Technology for Engineers,
https://doi.org/10.1007/978-3-031-32084-2_1
2 1 Engineering Materials, Atomic Structure and Bounding
and hence fall in the category of condensed matter. Gases are characterised by atoms/
molecules separated from each other by large distances.
There are some interesting consequences of the above-mentioned classification:
glass and several types of plastics like polyvinylchloride (PVC), for example, are
defined as rigid, supercooled liquids. Most materials of engineering interest are either
solids or rigid supercooled liquids.
Material science also classifies matter in four broad classes: metals, ceramics,
polymers and composites. Metals are characterised by their lustre, good conductors
of heat and electricity and to some extent by their property of ductility. A ceramic
is a material that is neither metal nor organic. It may be crystalline or glassy (rigid
supercooled liquid) or both. Ceramic pottery is quite well known but clay, bricks,
tiles, glass, concrete and cement are some other examples of ceramics. Ceramics,
depending on their composition may be semiconducting, superconducting, ferroelec-
tric or insulator; hence they are finding ever-increasing applications in solid state elec-
tronics, fiber optics, artificial joints, space shuttle tiles, micropositioners, chemical
sensors, body armours, self-lubricating bearings, etc. Polymers are mostly organic
substances made of long chains of molecules. Skin, hair and wood are examples of
polymers.
Another class of materials is called ‘composites’ that are combinations of two
or more of the above-mentioned metals, ceramics and polymers. Composites are
materials designed for specific goals to achieve a combination of properties not
found in any single material. Then there are advanced materials that are finding
applications in highly sophisticated technical fields like, electronics, space tech-
nology, computers, etc. Advanced materials include semiconductors, nanoengineered
materials and biomaterials, etc.
1.1.1 Metals
All materials are made up of atoms, either of the same or of different elements. Atoms
of different elements are characterised by their Atomic Number Z and Atomic Mass
Number A. Russian scientist Dmitri Mendeleev developed a table, called periodic
table, where elements were arranged in order of increasing atomic number from left
to right and from top to bottom.
Elements in periodic table are arranged in groups and rows such that elements
falling in a group exhibit similar chemical behaviour. Based on the observed simi-
larity in their chemical properties, elements have been grouped together as (i) alkali
metals, (ii) alkaline earth metals, (iii) transition metals, (iv) other metals, (v) halo-
gens, (vi) noble gases, (vii) rare earth and lanthanoid elements, (viii) non-metals and
(ix) actinoid elements. These different groups of elements are shown with different
colours in Fig. 1.1 that shows periodic table.
Materials of metal group are composed of one or more metallic elements, like
gold, titanium, nickel, copper, iron and aluminium and often also contain very few
atoms of non-metals like carbon, oxygen and nitrogen. These non-metallic atoms
1.1 Classification of Condensed Matter 3
Fig. 1.1 Periodic table of elements
are either deliberately mixed in controlled amount or are present as impurity in a

metallic crystal. Impurity atoms play a crucial role in altering the properties of the
metallic crystal. Atoms in metals and their alloys are arranged in a very regular
fashion. Metals, in general and in comparison to ceramic and polymers are dense,
have higher density. Table 1.1 shows the density of some typical metals, polymers,
ceramics and composite materials. Composite shown in the last column of the table
is made by mixing E-glass fiber of density 2.56 × 103 kg/m3 with cast polyamide
of density 1.1 × 103 kg/m3 in different ratios, and the product density varied from
(1.15 to 1.73) × 103 kg/m3 .
Some important characteristic properties of metals are (a) ductility, (b)
malleability, (c) lustre (d) large values of electrical and thermal conductivities and
(e) high melting and boiling points.
(a) Ductility is a mechanical property that may be described as material’s
amenability to drawing into wire. In material science, ductility is defined by
the degree to which a material can sustain plastic deformation under tensile
stress before fracture or failure. Ductility is an important quality from the point
of view of manufacturing, defining the suitability of the material for manufac-
turing operations such as cold working. It also tells how far the material can
absorb mechanical overload. Ductility is often measured through the percent
elongation and reduction in area at fracture in a tensile test. The fracture strain
Table 1.1 Density of different materials

Metal Polymer Ceramics Composite
Name of the Density Name of the Density Name of Density Name of Density in
material in kg/ material in kg/ the in kg/ the kg/m3 ×
m3 × m3 × material m3 × material 103
103 103 103
Osmium 22.59 Low-density 0.92 Boron 2.50
polyethylene carbide
Platinum 21.50
Aluminium 2.80 Cellulose 1.36 Sintered 3.00 Cast 1.15–1.73
diacetate silicon polyamide
nitride Plus
E-glass
fiber
Brass 8.50
Copper 8.96 Transparent 1.08 Zirconias 5.50
acrylonitrile
Germanium 5.30
Indium 22.50 Bark, wood 0.24
Iron 7.87
Lead 11.30
is defined as the strain at which a test specimen fractures during a uni-axial

tensile test. In uni-axial tensile test a bar of the specimen is pulled axially along
the length. As a result the length of the specimen increases while its area of cross
section decreases. At some value of the tensile deforming force, the specimen
undergoes fracture. Percentage elongation or engineering strain may be defined
as;
final length − initial length

Percentage elongation = × 100 (1.1)
initial length
The percentage reduction in area is given as;
Percentage reduction in area

Initial area of cross section − final area of cross section
= × 100 (1.2)
Initial area of cross section
Materials can undergo two types of fractures under tensile stress: brittle frac-
ture and ductile fracture. In brittle fracture the fractured ends have irregular
shape. The two types of fractures are shown in Fig. 1.2. Metals under tensile
stress undergo ductile fracture, as they have the property to withstand plastic
deformation. Metals have the ability to absorb more energy prior to fracture.
It is important to note that the property of ductility depends on the tempera-

ture at which tensile stress is applied. At some temperature, it is possible that a
ductile material may change from ductile to brittle, or vice versa; it is therefore
important to know the value of this critical temperature. In most cases of metals
it has been observed that at lower temperature they are less ductile or more
brittle, while their power of absorbing more tensile stress energy increases with
temperature and their ductile strength, therefore, increases with temperature.
The minimum temperature at which the metal changes from brittle to ductile is
known as ductile to brittle transition temperature or (DBTT).
(b) Malleability is a physical property of metals that defines their ability to be
hammered, pressed or rolled into thin sheets without breaking. This property
may also be defined as property of metals to deform and take new shapes
under compression. The malleability of a metal can be measured by how
much compressive stress or pressure the metal specimen can withstand without
breaking. Both the properties of malleability and ductility originate from the
crystal structure of metals, which in turn depends on the type of bonding between
atoms and molecules of the materials. Most of metals and their alloys have one
of the three types of crystalline structures (i) body-centred cubic (bcc), (ii) face-
centred cubic (fcc) or (iii) hexagonal close packed (hcp). Unit cells of these
structures are shown in Fig. 1.3.
On application of some stress the layer of atoms in metal slides over the layer
below without damaging the crystal structure, and the specimen assumes a new
shape. As shown in Fig. 1.4i, sliding is relatively easy in case of (hcp) and (bcc)
crystal structures as compared to the (fcc) structure (Fig. 1.4ii).
Fig. 1.2 Ductile and brittle

fractures
Fig. 1.3 Unit cells of a hexagonal close packed (hcp), b face-centred cubic (FCC) and c body-
centred cubic (bcc) crystal structures
Fig. 1.4 In case of metal

layer of atoms contained in a
plane may slide over the
layer below it
Sliding of atomic layers is possible only in metals because in metals and their
alloys atoms are bound through metallic bonding. In metallic bonding all valence
electrons of each atom get detached from its parent individual atom forming an
electron cloud around all positively charged metallic ions. Since electrons are
no more attached or associated with a particular atom, these electrons are called
delocalised electrons or free electrons. Since electrons in this electron cloud
are not associated with a particular positive atomic ion, the binding of indi-
vidual atomic ion with electron cloud is quite weak. Therefore, on application
of stress, layers of atomic positive ions may move relative to the layers below
or above it, without damaging crystal structure. With the change in the shape
of the metal specimen, the delocalised electron cloud also assumes a new shape
and orientation such that the overall structure, i.e. the crystalline structure of
specimen does not break.
Bonding between atoms of a material depends essentially on two factors; the size
of atoms and their relative separation. There are several types of bonding that are
found in different materials. We shall study more details about bonding in the
next section; however, for present it may suffice that in metals and their alloys
a special type of bonding called metallic bonding is found. To understand
the nature of metallic bonding let us start with two atoms of the same metal
with atomic number Z. When these two atoms are far apart, each atom has the
positive nucleus that is surrounded by a number of electrons such that the total
positive charge of the nucleus is counterbalanced by the total negative charge
of electrons; therefore, there must be Z electrons in each of the two atoms.
These atomic electrons are distributed around the nucleus at different distances
in different discrete energy states. The group of electrons farthest from the
nucleus is termed as valence electrons. Since they are farthest from the nucleus,
valence electrons experience a very weak force of attraction by the positively
charged nucleus, in scientific language one says that valence electrons in an
atom are least bound. If N is the number of valence electrons, then (Z − N)
electrons will be tightly bound with the nucleus and one defines the core of the
atom as the nucleus plus (Z − N) tightly bound electrons. Obviously, the net
charge on atomic core will be (+ Ne), where e is the unit of charge. Atomic
core, therefore, behaves as a positive ion as shown in Fig. 1.5a. For example,
the number of valence electrons in aluminium atom is 3, hence each positive
core or ion of Aluminium has a charge + 3e (3 × 1.6 × 10–19 C).
Fig. 1.5 a Structure of the core and valence electrons in an isolated atom b arrangement of positive
ions (cores) and delocalised electrons in a metallic specimen
Next let us consider that another similar atom is brought near to the first, so close
that the two atoms start feeling the presence of the coulomb fields of each other.
Since the cores of two atoms are tightly bound, they will not be much affected
by the presence of the other atom, but loosely bound valence electrons of both
atoms will feel almost equal force of attraction by the core of its parent atom as
well as the core of the other atom. Thus valence electrons will now get associated
to both atomic cores. This results in delocalisation of valence electrons, which
means that valence electrons (numbering 2N) are now not confined to the field of
one atom but are relatively free to move from the field of one atom to the field of
the other atom. As a result a bond is formed between the cores of the two atoms
making the system a diatomic molecule. Valence electrons now move around
the diatomic molecule in fixed orbitals. If a third atom is now brought very near
to the diatomic molecule, a tri-atomic molecule is formed having 3N number of
delocalised electrons circulating around three atomic cores in specified orbitals.
In this way a piece of a metal may be considered as a multiatom molecule with
delocalised electron cloud around it, see Fig. 1.5b.
Delocalised electrons may freely move from one core to the other and then to the
other, free to move within the electron cloud but they cannot leave the electron
cloud. It is because if any electron tries to leave the multiatomic molecule
electron cloud, the molecule develops a net positive charge and pulls the electron
back. Often one uses the term ‘free electrons’ for delocalised electrons as
they are not attached firmly to any atom of the specimen and are free to move
from the electric field of one atom to the electric field of any other atom of
the multiatom molecule, yet they are bound to the molecule. The delocalised
electrons are very large in number, but they move in some well-defined orbitals.
Each electron has a fixed discrete value of energy. Not more than two electrons
can have same energy as per Pauli’s exclusion principle. The energy differences
are very small and, therefore, electron energy levels are closed spaced as shown
in Fig. 1.6. As shown in this figure, energy levels up to some energy are filled with
electrons, but there are large number of empty levels. When delocalised electrons
absorb energy from some external source, for example if light is made to fall
on metallic surface delocalised electrons may absorb incident light photons and
may shift to their next higher excited state. Similarly, if the metal specimen is
heated delocalised electrons absorb energy and shift to higher excited states.
Availability of large number of empty levels for electrons plays important role
in metals.
The crystal structure of metals consists of regular arrangement of atomic cores
(or positive ions left after losing valence electrons) in two-dimensional arrays,
called crystal lattice, stacked one over the other in three dimensions. A rough
and enlarged version of metallic crystal structure is shown in Fig. 1.7.
Two lattices marked A and B are shown in the figure with large number of
delocalised electrons moving all around constituting electron cloud. Though
Fig. 1.6 Closely spaced

energy levels of delocalised
electrons
Fig. 1.7 A rough and

enlarged view of metallic
crystal
the lattices A and B are shown quite well apart in the figure, but in actual case
successive lattices almost touch each other. In metallic crystals positive ions in
a lattice are very strongly bound with each other; therefore, it requires large
amount of energy to break lattice structure. The strong binding of positive ion
cores in metallic crystals is provided by the cloud of delocalised electrons which
serves as a glue. However, the binding is much weak between two adjacent
lattices. That is why application of small stress, capable of overcoming bonding
between lattices, may slide lattices one over the other, without damaging lattices.
Sliding of lattices with respect to each other results in giving a new shape to the
specimen. Since cloud of delocalised electrons is firmly bound with lattices, it
readjusts its orientation and other parameter to give stability to the new shape.
It may thus be observed that properties of ductility and malleability in metals
originate from metallic bonding that is characterised by delocalised electron
cloud, large number of empty electronic states and very strong binding between
positive ions in crystal lattices.
(c) Different types of metal strengths
Strength of a material is judged by the ability of the material to resist deformation

and failure under the action of external forces. Some important and frequently
used strength indexes for metals are:
(i) Tensile strength It is defined as the maximum load that a metal part can
support without fracture when being stretched, divided by original cross
sectional area of the material. Mostly it is expressed in units of pounds per
square inch (PSI) or Pascal denoted as Pa. The unit, named after Blaise
Pascal, is defined as one newton per square metre.
(ii) Yield strength It is defined as the maximum stress a material component
can withstand without permanent deformation or the stress at the yield
point at which the specimen starts plastic deformation.
(iii) Compressive strength In contrast to the tensile strength, the compressive
strength of a material is defined as the maximum compressive stress a
solid material can sustain without fracture under gradually applied load.
It is a measure of the capacity of the material to withstand loads tending
to reduce its size.
(iv) Impact strength It is defined as the maximum impact or suddenly applied
force a specimen can take before its failure. As a matter of fact it is a
measure of how much energy the specimen may absorb at the limited
state.
(v) Shear strength It is defined as the maximum shear load a material can
withstand before failing divided by its cross sectional area.
(vi) Ultimate strength It is measured as the amount of utmost tensile,
compressive or shearing stress that a given unit area of the given specimen
sample can bear without deformation or breaking.
Table 1.2 shows the values of tensile and yield strengths for some metals.
(d) Lustre Freshly cut surfaces of all metals shine when light falls on them. This
property of metals is called lustre. This happens because of the delocalised
electron cloud and the availability of large number of empty electron energy
levels. Although the assembly of delocalised electrons is called a cloud but
electrons in this cloud move in systematic way. In electron cloud there are large
number of electrons, each having a discrete energy and moving in a defined
orbital. When light from some external source falls on these electrons, they
absorb incident light of all frequencies and get shifted to their next excited
level which is empty. However, electron excited states are very short-lived;
Table 1.2 Tensile and yield

Type of metal Yield strength in PSI Tensile strength in
strength for some metals
PSI
Aluminium-3003 21,000 22,000
Copper 28,000 –
Stainless steel-304 40,000 90,000
Titanium 37,000 63,000
excited delocalised electrons almost immediately revert back to their ground

states emitting the photons of nearly the same energy or frequency that they have
absorbed. As a result of the emission of light photons from the de-excitation
of electrons, metal surface shines. Surfaces of some metals like silver, gold,
etc. remain shining all the time as these metals do not chemically react with
environmental chemicals and gases. But in case of some other metals surfaces
shine only when they are freshly cut, after some time the metal chemically reacts
with atmospheric chemicals and a thin layer of metal oxide, etc. get deposited
on the surface making it dull.
(e) Electrical and thermal conductivities Metals are good conductors of electricity
and heat; they have large values for electric and thermal conductivities. Let us
consider a specimen of some material of length L and area of cross section A,
if now a voltage V is applied across the two opposite faces of the specimen
an electric current I may flow through the specimen. According to Ohm’s law,
current I will be proportional to V, i.e.
V ∝ I or V = R I (1.3)
In Eq. (1.3), the proportionality constant R is called the resistance of the given
piece of specimen and is a measure of the opposition that the specimen has
offered to the flow of current through it. A large value of R is an indication that
the specimen has offered large opposition to the flow of current, i.e. it is not a
good conductor of electricity. Resistance R is measured in units of Ohm (Ω). The
magnitude of R depends on (i) the material and on two physical parameters, (ii)
area of cross section A of the specimen and its (iii) length L. One may therefore
write,
L L
R∝ or R = ρ (1.4)
A A
The constant of proportionality ρ in above equation is called the resistivity

or specific resistance of the material and is measured in units of (Resistance-
metre) or (Ohm meter) also represented as (Ω-m). In Eq. (1.4) if one puts L =
1 m and A = 1 m2 , then R = ρ, that means that resistivity of the material is
equal in magnitude to the resistance offered between the two opposite faces of
a cube of side 1 m. Resistivity is a property of the material of the specimen,
and it does not depend on the physical size of the specimen. The reciprocal
of resistance R is called conductance and is measured in units of Siemen
denoted by S. 1 S = 1 ohm1 (Ω) . In older literature ‘moh’ has been used in place
of Siemen as unit of conductance. The reciprocal of resistivity, ρ1 , is called the
specific conductivity or simply conductivity of the material and is denoted
by σ . Conductivity is measured in units of (Ω-m)−1 which may be written as
Siemen per metre represented by (S/m).
Figure 1.8 shows the bar graph for the range of values of conductivity for
different materials. The conductivity of metals and alloys has high value but
varies in a narrow range. Conductivity has largest range for composites, as
expected. Conductivity value at room temperature (≈ 20 ◦ C) for some important
metals and alloys is given in Table 1.3. It may be observed in this table that
silver that has the maximum value of conductivity is the best conductor of
electricity and, therefore, it is frequently used in making electrical connections
in sophisticated electronic circuits.
Large number of delocalised or free electrons in metals is responsible for their
high electrical and thermal conductivities. Electron, being negatively charged,
experiences a force F = −eE when subjected to an electric field of strength
E. The negative sign in this expression tells that the force F is in a direction
opposite to the direction of electric field E.
Figure 1.9 shows a rectangular metallic rod of length ‘d’ which is connected
to a battery of voltage V with a switch ‘sw’. When the switch is made on, face
F 1 of the rod is connected to the positive terminal of the battery and face F 2
to the negative. A current of magnitude I flows through the metallic rod such
that V = R I . As already mentioned, R is a measure of the opposition that the
metallic rod offers to the flow of current. In the following we shall discuss the
mechanism of current flow and the origin of resistance R.
Fig. 1.8 Bar graph showing

range of conductivity for
different materials
Table 1.3 Conductivity for

Material Conductivity (S/ Material Conductivity (S/
some metals and alloys at
m) m)
20 °C
Silver 6.30 × 107 Manganin 2.07 × 106
Copper 5.96 × 107 Constantan 2.04 × 106
Gold 4.1 × 107 Nichrome 9.09 × 105
Aluminium 3.5 × 107 GaAs 5 × 10–8 to 103
Fig. 1.9 Current flow

through a metallic specimen
The interior of the rod has a crystalline structure consisting of positively charged
core of metallic atoms (indicated by + symbol) arranged in a regular fashion
in three dimensions. These sheets of positive ions are called crystal lattices and
they are held at their positions by the balance of attractive forces of delocalised
electron cloud and repulsive forces between nearby positive ions. Almost free
(delocalised) electrons, in motion, each with its inherent velocity (indicated by
brown coloured arrows) surround the crystal lattices. With the application of
voltage V across the two opposite faces of the rod electric field E of magnitude
E = Vd , directed from face F 1 to face F 2 , gets established between the two faces
F 1 and F 2 inside the rod. The electric field applies a force F = |eE| directed
from F 2 to F 1 , on each of the delocalised electron and imparts an additional
velocity, say vad directed opposite to the field direction, to each electron. This
additional velocity component on each electron is represented by a small black
arrow in Fig. 1.9. Under the action of two velocities, each delocalised electron
moves in the direction of the resultant velocity. However, the additional velocity
component tries to make every electron in the road to rush towards face F 1 . At
first it appears that all electrons in the metal rod will reach face F 1 in no time and
will accumulate there. But that is not true. These free electrons mostly moving
towards face F 1 collide with crystal lattice and their direction of motion and
the magnitude of velocity both get changed at such collisions. Free electron–
lattice collisions are very frequent, as a result though there is a net flow of
negative charge in the rod from face F 2 towards face F 1 but on average number
of free electrons per unit volume of the rod remains almost constant; there is no
accumulation of electrons in any part of the rod. The net flow of negative charge
(from F 2 to F 1 ) establishes the current I, and conventional current is assumed to
flow from F 1 to F 2 . One may ask a question why collisions between electrons are
not taken into account. The answer is that collisions between electrons are quite
unlikely because of their negligibly small size and very small time that electrons
take in crossing each other. Such electron collisions with crystal lattice are the
major cause for randomisation of electron velocities. Larger the frequency of
electron–lattice collisions, lesser will be the amount of net charge flow towards
face F 1 , resulting in lower current. It may, therefore, be realised that in metals the
opposition R to the flow of current originates from electron–lattice collisions.
Since inherent speed of electrons increases with temperature, the electron–lattice
collision frequency strongly depends on the temperature of the metallic rod; at
higher temperature the resistance R of the same specimen rod will be larger.
This is confirmed from the experimental observed fact that the same specimen
shows a larger value of resistance R or of resistivity ρ and a lower value for
conductivity σ at higher temperatures. It is, therefore, required to specify the
temperature while mentioning the resistance, resistivity or conductivity of a
specimen.
Large number of electrons colliding with lattice impart energy and momentum to
it that sets the crystal lattice in vibratory motion. This vibratory lattice motion
is quantised, i.e. the lattice in vibratory motion either absorb or emit energy
in packets. Energy packets corresponding to the vibratory lattice motion are
called ‘phonon’. There should be no confusion between phonons and photons;
photon is the energy quanta of electromagnetic field, while phonon is the quanta
of lattice vibration.
Heat is a form of energy, and it may transmit from one place to another by
three distinct methods; (i) radiation (ii) convection and (iii) conduction. Heat
transfer through radiation does not require any medium; it is directly transferred
as energy quanta, photons, from the source to the receiver. Sun light reaches
Earth via radiation crossing a vast region of vacuum. However, both convection
and conduction require some material medium to transfer heat. In convection,
medium particles take heat energy and move away to transport heat energy.
This happens in boiling of water when water molecules absorb heat from the
hot source at the bottom of the container and move out to transfer heat energy to
the top layer. In case of conduction on the other hand, medium particles do not
move, instead they transfer heat energy to the next particle, and then to the next
and so on. Though heat transfer from a hot body by radiations cannot be avoided,
but in solid materials heat transfer essentially takes place via conduction.
Thermal conductivity of a material, generally denoted by κ, tells about the
ability of the material to let heat energy pass through it via the process of
conduction. A large value of κ for a given material means that the material is
a good conductor of heat. Metals and their alloys are good conductors both of
heat and electricity.
Let us consider a rectangular sheet of a material of thickness ∆x (m) and area
of cross section A (m2 ) as shown in Fig. 1.10. Let the front face of the sheet
be at temperature (T + ∆T ) (Kelvin K) and the back face at temperature T
(K). Front surface being at a higher temperature will conduct some heat energy
∆Q (Joule J) in time ∆t towards the back surface. Experimentally it has been
found that heat transfer through conduction ∆Q from front face to the back
face is proportional to the area of cross section A of the surface, time ∆t, the
Fig. 1.10 Heat flow across

the thickness of a rectangular
sheet
temperature difference ∆T and is inversely proportional to the thickness ∆x.

One may, therefore write,
A . (∆t) . (∆T ) A.(∆t).(∆T )

−∆Q ∝ Or − ∆Q = κ .
∆x ∆x
Or
∆Q
κ =−
A . (∆t)(∆T /∆x)
× [Joule per unit area per unit time per unit temperature gradient]
× [whatt/metre Kelvin = W/m K] (1.5)
The left side of Eq. (1.5) gives the amount of heat energy lost by the front
surface per unit time per unit area, and, therefore there is a negative sign, the
negative sign simply shows that energy is lost by the front surface. The units
of thermal conductivity κ may be given as Joule per second per unit area per
unit temperature gradient or watt per meter per Kelvin, i.e. (W/m K). In metals
heat conduction also takes place through the delocalised electrons which are
relatively free to move and transport heat. Since the number of free electrons
in metals is large, the conductivity of metals and alloys is also large. Values of
thermal conductivity for some metals and other materials are shown in Table
1.4.
(f) High melting point of metals
Metals in general have high melting and boiling points. When a solid melts,
its crystalline structure gets destroyed and it changes its phase from solid to
liquid. On further heating the liquid, at some temperature called boiling point,
changes into gaseous phase. A high melting point means that the bonding
between constituent atoms is very strong. In case of metals crystalline structure
is protected by very strong bonding of multiatomic molecules, i.e. by the large
Table 1.4 Thermal conductivities of some materials

Material Thermal conductivity (W/mK) Material Thermal conductivity (W/mK)
silver
Silver 403 Iron 94
Copper 401 Nickel 106
Gold 327 Diamond 1000
Aluminium 237 Fiber glass 0.04
binding energy of crystal lattices. Tightly bound lattices in metallic crystals

require large amount of heat energy to break them that is why the melting and
boiling points of metals are high. Further, the binding energy will be larger for
the lattice that contains larger number of ions per mol. Hence, heavier metals
have higher melting and boiling points. Melting and boiling points of some
metals are shown in Table 1.5, which clearly shows the increase in both melting
and boiling points with the atomic weight of the metal.
When a piece of metal is looked through a high resolution microscope it is
generally observed that it does not contain a single large crystal. Instead, there is
large number of small crystals of different sizes packed together with different
orientations. These small crystals are called grains and the two-dimensional
surfaces between adjacent grains as grain boundaries. Grains have different
crystallographic orientations. Process of grain formation in a metallic specimen
starts when molten metal cools and crystallisation takes place. Small crystals
having different orientations start growing simultaneously at different locations
in the specimen and grow till they fuse. Grain size is an important parameter
for given metal. Small grain size increases tensile strength and tends to increase
ductility. However, grains of small size reduce the electric and thermal conduc-
tivities. Grain boundaries may be looked as 2-D crystal defects and tend to
reduce thermal and electrical conductivities. In general, grain orientations are
random, but they may be aligned to some extent by repeated rolling along one
direction.
SAQ: On heating metals mostly become more malleable, why?
SAQ: Ductility of metals that has grains of small size is more, how can this be
explained?
Table 1.5 Melting and boiling points of some metals

Material Melting point Boiling point Material Melting point Boiling point
(°C) (°C) (°C) (°C)
Aluminium 660 2515 Iron 1150–1593 2861
Silver 961 2162 Platinum 1770 3825
Copper 1084 2562 Tungsten 3400 5550
SAQ: Crystal lattice consists of positively charged ions but they are very tightly
bound. Which forces provide this strong binding?
SAQ: Current in metals is constituted by the flow of electrons; what will happen to
these electrons when a current carrying wire is cut, will electrons go out of
the wire?
SAQ: How can one explain the reduction of conductivity in metals that have small
grains?
SAQ: Why do delocalised electrons in metals remain inside the specimen when
they are not attached to any specific atom?
1.1.2 Ceramics
It is difficult to define ceramics in the present context as ceramics cover a very

wide range of inorganic materials that may contain metallic or non-metallic chem-
ical elements and are produced by many different physical and chemical processes.
There was a time when ceramics were defined as non-organic, non-metallic materials,
having very low electric conductivity and produced essentially by high-temperature
treatment. Some old literature also tried to define ceramics as ‘refractory’ mate-
rials, which in technical material science language means, materials capable of with-
standing every day abuses like extreme temperatures, attacks from acid and alkalis
and general wear and tear. However, the old definition is no more valid, today ceramic
materials having metallic ions and novel properties including semiconductor and
superconductivity produced by many different techniques are being used in indus-
trial applications. Some broad features of ceramics, as compared to polymers and
metals, are shown in Table 1.6.
Sometime it becomes easier to define materials in terms of their properties; their
behaviour on heating, on passing current through them, or putting them in water,
etc. Such a classification becomes confusing, for example graphite, an allotrope of
carbon, is considered a ceramic because it is non-metallic and inorganic. However,
unlike most ceramics it is soft, wears easily and is a good conductor of electricity.
Diamond which is another form of carbon on the same grounds is treated as ceramic,
but properties of diamond are totally opposite of graphite; it is hard, very stable and
does not wear out easily. With regard to mechanical behaviour, ceramics are relatively
stiff and strong; stiffness and strength are comparable to metals. In addition, ceramics
Table 1.6 Broad features of Ceramics

Material Property
Chemical stability Density Melting point (°C) Plasticity
Ceramic Non-reactive Low to medium Medium to high Brittle
Polymer Very reactive Very low Low Ductile to brittle
Metal Reactive Medium to high Low to high Ductile
are very hard, extremely brittle and susceptible to fracture. Figure 1.11 shows the
spread of tensile strengths for different materials. In spite of being very hard, ceramics
may be optically transparent, translucent or opaque.
Type of bonding between constituent atoms is often used to classify a material; for
example, metals are characterised by metallic bonding where delocalised electrons
of constituent atoms provide the ‘glue’ for strong binding. Similarly, polymers have
strong covalent bonds in long molecular chains, while relatively weak van der Waals
bond binds one long molecular chain with the other. On the other hand, ionic bonds
are found in non-metals. In ceramics all types of bonds exist. As an example, Al2 O3 ,
MgO, SiO2 , etc. have ionic bonds, while SiC, BiC, BN, Si3 N, Si2 N2 O, etc. show
covalent bonding. It is interesting to observe how relative content of different types
of bonds changes the melting point of ceramics (see Table 1.7).
Solids generally have crystalline structure with grains, i.e. they are poly crys-
talline; microscopically there are several crystals having different orientations, fused
together. Fine-grained pure alumina and glass are polycrystalline ceramics. Ruby,
diamond, etc. on the other hand are large single crystal ceramics.
Ceramics have many different chemical compositions:
There are simple oxides that have high melting point, like ThO2 (melting point
Tm = 3300 °C), MgO (Tm: 2825 °C), UO2 (Tm: 2810 °C), etc. Ferrites that are
complex oxides like Fe3 O4 , SrF12 O19 ; Titanates: BaTiO3 (Tm: 1625 °C); SrTiO3
(Tm: 2080 °C); Nitrides: Si3 N4 (Tm: 1900–2600 °C); TaN (Tm: 3080 °C); Brides:
Fig. 1.11 Bar graph

showing the spread of tensile
strengths of different
materials
Table 1.7 Effect of bond

Ceramic Ionic bond Covalent bond Melting point
type on melting point of
compound (%) (%) (°C)
ceramics
SiC 11 89 2830
Si3 N4 30 70 1900
SiO2 51 49 1715
HfB2 (Tm: 3350 °C); ZrB2 (Tm: 3245 °C); Silicides: Hf5 Si3 (Tm: 2600 °C); WSi2
(Tm: 2160 °C); Halides: NaCl (Tm: 800 °C); Intermetallides: HfRe2 (Tm: 3160 °C),
Nb3 Sn; Metal ceramics: WC-TiC-Co; Polymer ceramic: Synthetic resins.
Material scientists occasionally divide ceramics into Traditional and Advanced.
Bricks, pottery, glass, porcelain, etc. are time tested daily use general purpose
ceramics. Out of these pottery is generally made from traditional clay while bricks,
tiles, etc. are heavy clay products. However, coarse-grained refractors fired bricks;
silica bricks find special use in making high-temperature oven, etc. Cement and
concrete are other traditional ceramics used in building construction.
Advanced ceramics are those that have been specially engineered, mostly since the
early twentieth century, for highly technical and specific applications. For example,
aluminium nitride (AlN) and beryllium oxide (BeO) ceramics have been devel-
oped to serve as heat-sink for electronic elements. These ceramics are also used as
substrate in electronic packaging. Similarly, SiO2 and polymer-ceramic compounds
are used as thermal protection shields. High-resistivity conductors-ceramics silicon
carbide (SiC), zirconium oxide (ZrO2 ), molybdenum silicate (MoSi2 ), etc. are used
for making heating elements and electrodes.
Ceramics having magnetic and superconducting properties have also been devel-
oped. Complex ferrite and oxides of heavy metals like, (Ba, Sr) Fe12 O19 ; Y Co5 Sm2
(Co, Fe, Cu, Zr)17 and Nd2 Fe14 B, etc. are used for making hard magnets.
Magnet-doped silicon dioxide (SiO2 ) and chromium-doped complex Be3 Al2
(Si3 )6 have been used to make artificial gem stones, former as topaz and the later as
emerald. Artificial diamonds are made from ceramic ZiSiO4 .
Advanced ceramics are finding extensive use in nuclear technology. UO2 , UC and
PuO2 are used as nuclear fuels. Ceramics BeO, BeC2 and ZrO2 have been used to
moderate fast neutrons to thermal energies in nuclear reactors. Similarly, ceramics
B4 C, HfO2 and Sm2 O5 are used for neutron shielding.
Some ceramics are biocompatible that means they are not harmful or toxic for
living tissues. Such ceramics are used for making artificial joints, prostheses, cardiac
valves and other implants.
One may conclude by saying that ceramics are versatile materials that have
applications in almost all walks of human life.
1.1.3 Polymers
Polymers are long-chain molecules of very high, running into hundreds and
thousands, molecular weight. It is for this reason that they are also called
‘macromolecules’. In old literature, when polymer science was not so developed,
term ‘resins’ was used for polymers. As has been mentioned, polymers are substances
made up of recurring structural units, each of which is regarded as derived from a
specific compound. These building blocks or units are called monomer. The physical
and the chemical properties of a polymer depend very strongly on the number of
monomer units in the chain. As an example take a simple case of normal alkane
hydrocarbon series;
The bond structures of first three members of the series are shown in Fig. 1.12a,
the monomer of the series is given in figure (b), while the series formula is shown
in part (c) of the figure. In series formula ‘n’ gives the number of monomers in that
particular member. It is interesting to note that the physical and chemical properties
of different members of the series change with the number ‘n’ of monomers in the
chain. First four members of the series are gases. The fifth member ‘n-Pentane’ is a
low viscosity fluid with boiling point of ≈ 36 °C. With the increase of the number
of monomers in the chain, the viscosity and boiling points of the member increase.
Some characteristics of series members and the value of ‘n’ for them are tabulated
in Table 1.8.
The boiling point of successive members of the alkane chain also increases with
the number of monomers, but the rate of increase slows down such that the boiling
point for the massive members of the series saturates at about 145 °C. Long-chain
alkanes having 103 to 3 × 103 carbon atoms are known as polyethylene. A big
difference between wax and polyethylene lies in their mechanical behaviour. While
polyethylene is a tough plastic, wax is a brittle crystalline solid.
Fig. 1.12 a Bond structure

of first three members of
alkane series, b monomer of
the series, c general
representation for the series
Table 1.8 Variation in some

Number of monomers in the chain Physical properties
properties of alkane series
with number of monomers 1–4 Gases
5–8 Low-viscosity liquids
9–16 Medium-viscosity liquids
17–25 High-viscosity fluid
26–50 Crystalline solid, wax
50–1000 Semicrystalline solids
1000–5000 Tough plastic solid,
polyethylene
~ 105 Fibers
The difference in mechanical properties arises of their structures. Figure 1.13

shows that in case of wax linear carbon chains of up to 50 carbons are linked with
each other by van der Waals forces that are quite weak. Hence, wax is brittle. On the
other hand in polyethylene monomer chains are quite big and are folded, as shown in
Fig. 1.14. Further, several chains are linked together through entanglements which
are difficult to break.
The process of formation of a polymer from its monomer is called polymerisation.
The number of monomers in a given polymer chain ‘n’ is called the degree of
polymerisation (DP). The functionality of a monomer is the number of sites it
has for bonding to other monomers under the given conditions of polymerisation
reaction.
While the exact molecular weight required for a substance to be called a polymer
is a subject of continued debate, often polymer scientists put the number at about
25,000 g/mol. Polymers may be classified in many different ways; for example,
Fig. 1.13 Structure of wax
Fig. 1.14 Structure of

polyethylene
classification based on (i) molecular forces, (ii) heat treatment, (iii) source, (iv)
structure and (v) mode of polymerisation, etc.
(i) Classification based on molecular forces Two types of bonds are frequent in
polymers; the hydrogen bond and (b) van der Waals bond. These bonds bind
chains and monomers in a chain with each other. Accordingly, polymers may
be classified as:
(a) Elastomers Rubber—like solids fall in this category. In these polymers
chains are coupled with each other by the weakest intermolecular force
which permits the polymer to be stretched. However, there are some cross-
links between the chains that bring back the polymer to its original shape
when deforming force is withdrawn. Examples are neoprene, buna-N, etc.
(b) Fibers Fibers are those polymers which may be drawn into long filaments
with lengths at least 100 times of their radii. This happens because of
the strong bonds between chains, usually hydrogen bonds. As a result
of strong intermolecular force, these materials are closely packed and
have crystalline structure. Examples are polyesters (terylene), polyamides
(nylon), etc.
(c) Resins They are liquid polymers that are used as adhesives, sealants, etc.
Examples are epoxy adhesives and polysulphide sealants.
(ii) Classification based on heat treatment A polymer that may be given different
shapes to make tough and hard utility articles by heating and/or by applying
pressure is called plastic. Plastics may be further classified as: (d) thermoplastic
(e) thermosetting plastic.
(d) Thermoplastic polymers Some polymers become soft on heating and
can be given any desired shape. However, on cooling they again become
hard and tough. The process of heating, reshaping and becoming tough
and hard on cooling can be repeated several times. The intermolecular
forces in these plastics are stronger than that in elastomers and weaker
than those in fibers. Sealing wax, nylon, PVC, etc. are some examples.
(e) Thermosetting polymers Those plastic polymers that undergo some
chemical changes on heating and become infusible mass which cannot be
given any shape are called thermosetting plastic polymers. In such poly-
mers, heating creates large number of new cross-linking bonds that convert
it into an infusible mass. Bakelite is an example of such thermosetting
polymer.
(iii) Classification based on source Based on the source of the polymer there may
be three classes:
(f) Natural Polymers Polymers that are found in nature, in plants and animals
are called natural polymers. Proteins, cellulose, barks, starch, and rubber,
etc. are some examples.
(g) Semisynthetic polymers Derivative cellulose that is obtained by modi-

fying it is called semisynthetic polymer. Cellulose acetate also called rayon
and cellulose nitrate are examples of semisynthetic polymers.
(h) Synthetic polymers Man-made polymers, like bakelite, polythene,
synthetic rubber, etc. are examples of synthetic polymers.
(iv) Classification based on structure On the basis of their structure polymers
may be called:
(i) Linear polymers These polymers contain long and straight chains of
monomers with very little or no cross-linking of chains. Examples are
high-density polythene, PVC, etc.
(ii) Branched polymers These polymers have linear chains with few
branches; examples are low-density polythene.
(iii) Cross-linked polymers These polymers are made up of bi-functional
and tri-functional monomers and, therefore, have strong covalent bonding
between linear monomer chains. Cross-linked polymers are usually hard,
do not melt or soften on heating and mostly do not dissolve. Examples
are vulcanised rubber, formaldehyde resins, etc.
(v) Classification based on mode of polymerisation Polymers may also be
classified on their mode of polymerisation into two classes;
(l) Addition polymers This type of polymers is produced by repeated addi-
tion of monomer molecules that have double or triple bonds. Addition
polymers formed by the polymerisation of a single monomeric species
are called homopolymer; polythene is an example of the same.
nCH2 = CH2 →→ Polymerization →→ −(CH2 − CH2 )n −

Ethene Polythene (Homopolymer)
Polymers made by addition polymerisation of two different monomers

are called copolymers; examples are Buna-S, Buna-N, etc. Figure 1.15
shows the addition polymerisation of two different monomers butadiene
and styrene into Buna-s copolymer.
(m) Condensation polymer When two different bi-functional or tri-
functional monomers undergo repeated condensation reactions, they form
a condensation polymer. In this type of polymerisation small molecules
like, water, HCl, alcohol, etc. are usually eliminated. Examples are nylon
6.6, etc. If one of the two monomers is tri-functional or there are three
different bi-functional monomers and they undergo condensation poly-
merisation, the resultant polymer has linkage sequences in two or three
dimensions; they are called cross-linked polymer.
Fig. 1.15 An example of

copolymerisation
SAQ: Name three types of bonds that are found in polymers.

SAQ: What is the class of the polymer that converts into an infusible mass on
heating? What is the reason for this change?
SAQ: What is meant by the functionality of a monomer?
SAQ: What do you understand by cross-linking in polymers? How does it happen?
SAQ: What must be the properties of monomers that on polymerisation produce
linked polymers?
1.1.4 Composites
Modern technologies often require materials with very special properties which are
not available in metals, metal alloys, ceramics and polymers. For example, aerospace
scientists are always in lookout of materials which have very low density, very strong,
highly resistant to abrasion and impact yet quite stiff. This amounts to asking for
two apparently opposite characteristics in the same material, because strong and stiff
materials are generally dense. Further, increasing the strength or stiffness, in general,
decreases the impact strength. An answer to such problems comes from composites;
these are materials produced by the combination of two or more of the three materials,
namely metals, ceramics and polymers.
A composite may be defined as a combination of two or more materials (often
called phases) at a microscopic scale and have chemically distinct phases that results
in better properties than those of the individual components used alone. Though
heterogeneous at a microscopic scale, a composite is statistically homogeneous at
macroscopic scale. In general, out of different phases in a composite, one partic-
ular material has volume wise larger concentration than others. This component
with largest concentration (or bulk material) is called the ‘matrix’. The other mate-
rial which is in relatively smaller amount is termed as the ‘reinforcement’. Rein-
forcements are primarily added to increase the mechanical strength, toughness and
stiffness of the material.
The manmade composites may be divided into three categories, (a) polymer
matrix composites (PMC), (b) metal matrix composites (MMC) and (c) ceramic
matrix composites (CMC).
(a) Polymer matrix composites Some polymers, in particular epoxies and
polyesters, have a notable property that they may be easily moulded into desired
complex shapes. But their drawback is that they do not possess high mechanical
strength as metals. On the other hand materials like glass, boron and aramid have
extremely high tensile and compressive strengths which are, however, not readily
apparent in their solid forms. This happens because when stressed, randomly
distributed surface ‘flaws’ (abnormalities due to impurity, etc.) in these mate-
rials make the solid crack or break much below its theoretical ‘breaking point’
stress. To overcome this problem, fibers of these materials (boron, glass and
aramid) are drawn. The advantage in fibers is that though the random distri-
bution of faults will still be same but only few fibers will be affected by these
faults and a very large number of fibers will have no fault in them and will
break at their theoretical break point. Therefore, a bundle of fibers will reflect
more accurately the optimum performance of the material. It may, however, be
realised that a bundle of fibers will show its tensile strength only in the direction
of its length, just like in a rope. When these fibers are mixed as reinforcement
with a polymer matrix, like that of polyester, the resulting composite shows
exceptional mechanical strength comparable or even more than that of metals.
When a stress is applied to the composite, the matrix material, polyester in this
case, spreads the stress to fibers. Further the bulk matrix protects fibers from
atmospheric wear and tear, abrasion and impact. High strengths and stiffness,
ease of moulding into complex shapes makes the composite superior to metals in
many ways. The overall strength of the composite depends on following factors:
(i) Properties of the matrix polymer
(ii) Properties of the reinforcing fiber
(iii) The ratio of the fiber to the polymer, called fiber volume fraction (FVF)
(iv) Geometry and orientation of fibers in the matrix.
It is obvious that the mechanical strength of the composite will increase with
the increase of FVF, but there are limits to which this ratio may be increased.
Firstly, it is essential that all fibers must be fully rapped with polyester matrix
from all sides so that they are not exposed. Further, the manufacturing process
that involves mixing of reinforcement with matrix often produces faults and
air-inclusion, which may become cause of breakdown. In case of ordinary appli-
cations, like boat-building industry FVF of 30–40% is quite enough. However,
in more sophisticate applications like aero-industry FVF of around 70% have
been obtained by advanced manufacturing methods.
The orientation of the fiber in the composite is also important because the
maximum tensile strength of the fiber is along its length and the tensile strength
in direction normal to the length is negligible. The composite is, therefore,
anisotropic which is in contrast to metals and alloys which are largely isotropic.
It is, therefore, very important when considering the use of the composite at the
design stage to know the magnitude and direction of the load in the finished
structure. If properly taken into account, the property of anisotropy of compos-
ites may be used to advantage, as composite material may be used only where
there are locked stresses.
There are four main types of direct loads that a composite may have to bear in
a structure. They are;
(i) Tension Load Fig. 1.16a shows the situation when tensile force is applied
to a composite. Response of the composite to tensile load very much
depends on the tensile strength of fiber reinforcement mixed with the
polymer, since it is much higher than that of the matrix material.
(ii) Compressive Load Application of compressive load to a composite is
shown in Fig. 1.16b. In this case the adhesive and stiffness properties of
the matrix polymer are very important as they have to maintain the fiber
straight and to prevent them from buckling.
(iii) Shear Load As shown in Fig. 1.16c a shear load attempts to slide adjacent
layers of the reinforcement fiber over each other. In this case also the
properties of the matrix polymer plays a crucial role, it should not only
have good mechanical strength but should also have very good adhesive
force with fiber so that it remains firmly attached with it.
(iv) Flexure Load As shown in Fig. 1.16d, flexure load is a combination of
tensile, compression and shear loads. Therefore, both the matrix and the
reinforcement fiber must possess good adhesive and mechanical strengths.
(b) Metal matrix composites Metal matrix composites have found usage in our
lives from olden times. Metals like cast iron with graphite, steel with high
carbide contents are all examples of metal matrix composites. Artefacts made
of metal matrix composites as swords, body armours, chains, etc. are all found
in excavation of old habitation sites.
Fig. 1.16 Four types of loads that may be applied to a composite

There are many ways to classify metal matrix composites. One very often used
classification is based on the nature of reinforcement component; particles, layer,
fiber. Fiber composites may further be classified as, continuous fiber composite
and whisker composite materials. The continuous fiber metal composites may
either be monofilament or multifilament types.
The reinforcement material in metal matrix composites may have different
objectives. The reinforcement by light metals opens up the possibility of the
application of these light metal reinforced metallic composites in areas where
weight reduction is the first requirement. Frequently used light metals as rein-
forcement are Al2 O3 and SiC. The development objectives of light metal
reinforced composites are;
(i) Increase in yield strength and tensile strength at room temperature and
at higher temperatures, maintaining the minimum toughness or ductility
(ii) Increase fatigue strength particularly at higher temperatures
(iii) Increase in Young’s modulus
(iv) Reduction of thermal elongation
(v) Improvement in corrosion and thermal shock resistances
(vi) Low density
(vii) Mechanical compatibility with the matrix metal (thermal expansion
coefficient that matches with matrix metal)
(viii) Chemical compatibility
(ix) Good process ability
(x) Economic efficiency.
Some of the above-mentioned objectives may be achieved by using non-
metallic inorganic reinforcements. Ceramic particles, fibers and carbon fibers are
frequently used as reinforcement materials in metal matrix composites (MMC).
(c) Ceramic matrix composites (CMC) Ceramic matrix composites mostly
consist of ceramic fibers embedded in a ceramic matrix forming a ceramic
fiber reinforced ceramic composite (CFRC). Carbon and carbon fibers that are
also considered ceramic materials, along with fibers of other ceramic materials,
have been used as reinforcement elements. Typical reinforcing fiber materials
are; Carbon C, Silicon carbide SiC, Alumina Al2 O3 , Mullite or Alumina Silica
Al2 O3 –SiO2 .
Normally a ceramic matrix gets fractured by a tensile stress that produces an
elongation of about 0.05% in the length. However, if normal ceramic matrix is
reinforced by ceramic fibers, the fracture or cracks produced by excessive tensile
stress get covered up by the extension of fibers. An essential requirement for
complete recovery of the fracture site is that the matrix ceramic should also be
able to slid and fill the fracture gap. This requires that the adhesive force between
matrix and fiber is not very strong. A strong bond between the matrix and the
fiber will require a very high elongation capability of the fiber bridging the
fracture gap and would result in a brittle fracture. The adhesive force between
the matrix ceramic and the fiber is reduced by coating the fiber with a thin
Fig. 1.17 Ceramic matrix

composite with fiber
reinforcement
layer of pyrolytic carbon or boron nitride. These coatings weaken the bound
at the fiber–matrix interface. Figure 1.17 shows how a ceramic matrix fiber
reinforced composite repairs fracture or crack sites caused by tensile stress by
the elongation of fiber and sliding of the matrix ceramic.
Ceramic matrix reinforced by fibers composites may display both; the high
insulating or high conductivity properties. As a matter of fact the thermal and
electrical properties of ceramic matrix composites strongly depend on the prop-
erties of its constituents, namely fibers, matrix, pores in matrix, etc. Further,
fibers bring in anisotropy in behaviour of composites. Oxide ceramic matrix
composites are very good insulators. Because of their high porosity their thermal
insulation is much better.
The use of carbon fibers increases the electrical conductivity, provided the fibers
remain in contact with each other and with the voltage source. Silicon carbide
(which is a ceramic and a semiconductor) matrix is a good thermal conductor.
Electrical conductivity of SiC matrix decreases with the rise of temperature as
it is semiconductor.
Some important properties of ceramic matrix composites are;
(1) High thermal shock and creep resistance
(2) High temperature resistance
(3) Excellent resistance to corrosion, wear and aggressive chemicals
(4) High tensile and compressive strengths, thus no sudden failure as compared
to conventional ceramics.
Applications CMC have a wide range of applications, some of which are given
below;
• High-performance breaking systems
1.2 Atomic Structure 29
• Heat exchangers
• Bullet proof armour
• Turbine blades
• Heating elements
• Gas-fired burner parts
• Hot pressed dies
• Stator vanes
• Thrust control flaps for jet engines
• Refractory components
• Filters for hot liquids
• Heat shield systems for space vehicles
• Rocket propulsion components
• Turbo jet engine components.
SAQ: Fibers used as reinforcement in ceramic matrix composites are coated or
painted with some material. What is the need of this coating or painting?
SAQ: What is ‘meant by temperature shock’? Why CMC used in break lining should
have high resistance for temperature shock?
SAQ: What are ‘light metal MMC’? Which light metals are often used?
SAQ: Which polymers are frequently used as matrix material in PMC and why?
1.2 Atomic Structure
Though from engineer’s point of view it is only important to know different properties
of materials so that an appropriate specimen may be selected for the required use, but
if it is required to modify some property or to develop a new material having desired
properties, it is essential to know how and why different materials have different
properties. Key to this lies in the atomic and the molecular bonding of different
materials, i.e. how atoms and molecules are held together in different materials. Some
details of atomic structure along with different types of primary and secondary bonds
are discussed in the following.
1.2.1 Elements of Atomic Structure
All materials are made up of molecules, each molecule in turn, is made up of atoms.
Each atom when looked from a distance appears electrically neutral. However, on
a closer look, each atom has at its centre a nucleus with positive charge Ze. Here,
‘e’ stands for a unit of charge e = 1.6 × 10–19 C. Nucleus contains certain number
N of neutrons, each neutron being neutral, and Z number of proton each with + 1e
charge. Total number of nucleons (neutron and protons together are called nucleons)
in a nucleus is denoted by A, called atomic mass number and A = (N + Z). Number of
protons Z in a nucleus decides the amount of positive charge on the nucleus and (Z) is
called the atomic number of the nucleus/atom. The nucleus of the atom is surrounded
by a spherical or nearly spherical distribution of negatively charged cloud made of
electrons, each electron denoted by the symbol ‘e’, has − 1e unit of negative charge.
Symbol e is used to denote both the unit of charge as well as an electron, but this
does not create any confusion as the context of its use immediately tells whether it is
used for denoting electron or for charge. The total number of electrons in this cloud
is Z, so that the total negative charge surrounding the nucleus of atomic number Z
is – Ze.
When looked from a distance (much larger than the size of the atom), both the total
positive charge in the nucleus (+ Ze) and the total negative charge (− Ze) contained
in electron cloud appear as if they are held at a point at the centre of the atom (centre
of the nucleus). Since total negative charge is equal in magnitude to the total positive
charge, the net charge at atom’s centre becomes zero. Thus atom, looked from a
distance, appears electrically neutral.
1.2.2 Arrangement of Electrons in Atom
Quantum mechanical model of atom was proposed by an Austrian scientist named

Erwin Schrodinger in 1926. The model is based on formalism or mathematical recipe
that has some axioms or postulates. These postulates have no foundation in clas-
sical physics. Correctness of these postulates is derived from the correctness of the
predictions of the model. According to the quantum or wave model of the atom,
the presence/motion of electrons in the atom is completely described in terms of a
mathematical function called wavefunction of electron, denoted by Greek letter ψ.
The wavefunction is supposed to contain all information about the electron which
may be obtained by solving a differential equation called Schrodinger wave equation.
Wavefunction ψ(r ) is a function
( of the) distance ‘r’ of the electron from the centre
of the atom. The quantity ψ(r )∗ ψ(r ) gives the probability of finding an electron
at a distance ‘r’ from the centre of the atom. Here ψ ∗ is the complex conjugate of
ψ. Though in principle it is possible to write down and solve Schrodinger equa-
tion for any atom, however, it becomes complicated to do it for an atom that has
many electrons. Therefore, for most of the time one solves Schrodinger equation for
the simplest atom; the Hydrogen atom, the atom that has one proton in the nucleus
and one electron moving around it. Results obtained for Hydrogen atom are then
extended, with suitable modifications, for other atoms.
On solving Schrodinger equation for hydrogen atom one gets a number of
wavefunctions that are characterised by three quantum numbers, namely the prin-
cipal quantum number ‘n’, azimuthal quantum number ‘l’ and the magnetic
quantum number ‘ml ’. Each electron in the atom has a unique set of values for these
quantum numbers (n, l and m l ). A fourth quantum number called magnetic spin
quantum number (or simply spin quantum number) and denoted by m s is added
to the list of three quantum numbers obtained by solving Schrodinger equation. This
additional quantum number does not appear in the solution of Schrodinger equation
but is added to account for the two possible spin orientations of the electron. Electron
has an inherent spin of value 21 ℏ, here ℏ is quantum mechanical unit of measuring
spins. Magnetic spin quantum number (or simply spin quantum number) m s in case
of electron can have only two possible values; + 21 ℏ or − 21 ℏ. Although in principal
it is not possible to understand any quantum mechanical processes in classical terms,
however, for the sake of understanding, the two inherent spin motions of electron
may be associated with clockwise and anticlockwise directions of spin.
Microscopic systems or entities that follow quantum mechanics also obey a law
or principle called Pauli’s exclusion principle, according to which two electrons
in a given system (or atom) cannot have the same values for all the four quantum
numbers.
[ ∗ The region] of three dimensional space around the nucleus where the func-
tion ψn,l,m l
ψn,l,m l has maximum value is called the atomic orbital or simply
orbital. Orbital is a region of space around the nucleus of the atom where probability
of finding an electron with specified quantum numbers is a maximum. Classically,
orbital may be associated with classic orbit or Bohr orbit of the electron. But with the
difference that classic electron orbit is a well-defined sharp circular/elliptical path in
which electron travels around nucleus, while orbital is a volume of space around the
nucleus where the probability of finding electron with given set of quantum numbers
is maximum. Since there may be many different combinations of electron quantum
numbers, there are several orbitals for an atom.
Let us understand the physical significance of these quantum numbers.
(a) Principal quantum number ‘n’ Principal quantum number ‘n’ defines the
energy level of the electron or principle shell in atom. In quantum mechanics
particles can have only discrete values of energy. Principal quantum number ‘n’
can have only positive non-zero integer values, i.e. n may have values: 1, 2, 3,
4 and so on. Principal quantum number ‘n’ also determines the mean distance
of the electron from the centre of the atom that is from the nucleus. An energy
level with principle quantum number ‘n’ may accommodate a maximum of 2n2
electrons. Thus,
Energy level for which n = 1, may have at the most 2 electrons.
…
…
All electrons in a level of given principal quantum number ‘n’ have very nearly
same energy, but their other quantum numbers (l, m l , and m s ) are different.
As a matter of fact the maximum number of electrons 2n2 in energy level ‘n’ is
nothing but the number of different valid combinations of the remaining three
quantum numbers (l, m l , and m s ).
Let us now understand what is meant by the energy of an electron in an atom. In

an atom electrons are held because of the attractive force between the positively
charged nucleus and negatively charged electrons. This attractive force binds
the electron with the nucleus; it means that some amount of work will have to
be done (or some energy has to be spent) to take a given electron out of the grip
of the nucleus. An electron that is not bound with the nucleus is a free electron,
its binding energy with atom is zero. Electrons in an atom that are bound with
the nucleus have negative binding energies. The electron that is nearest to
the nucleus has the maximum negative (binding) energy or minimum absolute
energy. The electron (in an atom) which is farthest from the nucleus will have
minimum negative (binding) energy or maximum absolute energy. So when one
talks about the energy of an electron in the atom it means the binding energy of
the electron which is always negative. An electron near to the nucleus has less
absolute energy while the one far away has more absolute energy. Therefore,
absolute (binding) energy of electron in an atom increases from inner most to
the outer most electrons.
(b) Azimuthal quantum number l Azimuthal quantum number is related to the
shape of the orbital. It may have only positive integer values including zero up
to (n − 1). That is for a given value of ‘n’, l may have values; 0, 1, 2, 3, …,
(n − 1). For example, if n = 1, then l may have only one possible values 0. If
n = 2, then l may have two values 0 and 1, for n = 4, l = 0, 1, 2 and 3 and so
on. Each value of l defines a sub-shell or sub-state within the principal energy
level or shell or orbital defined by n.
Different numerical values of l are assigned different lower case alphabets;
l = 0; is called s-orbital; l = 1 is called p-orbital; l = 2, is d-orbital; l = 3,
f-orbital etc.
It is desirable at this stage to introduce the concept of multiplicity. As already
mentioned, the azimuthal quantum number l defines the shape of the orbital
in three-dimensional space around the nucleus. An orbital of a given shape
may have several different orientations with respect to the nucleus. In quantum
mechanics space is also quantised, and therefore, orientations are also quantised.
It is found that an orbital defined by azimuthal quantum number l can have
(2l + 1) different orientations. Which means that s-orbital can have only one
orientation as for s state l = 0 and, therefore, (2l+1) is 1. The p-orbital (l = 1)
may have (2 × 1 + 1 =) 3 different orientations and orbital f for which l = 3,
(2 × 3 + 1 =) 7 different orientations. The factor (2l + 1) which gives the
number of different orientations is also called the Multiplicity of the orbital.
(c) Magnetic quantum number ml This quantum number specifies the above-
mentioned different orientations of orbitals. For a given value of l the magnetic
quantum number m l may have (2l + 1) different values, starting from −l to + l
in steps of unity. For example if l = 3 (f-orbital), then quantum number m l
may have (seven different) values −3, −2, −1, 0, +1, +2, and + 3. Similarly,
if l = 1, which means p-orbital, the three different values of m l will be: − 1,
0, + 1.
(d) The Magnetic spin quantum number m s As already mentioned, m s does not
arise from Schrodinger equation, it is included to specify direction of the inherent
spin of the electron. In classical term, if the electron is spinning in clockwise
direction than ms = + 1/2 and if it is spinning in anticlockwise direction then
ms = − 1/2. These assignments of + 1/2 and − 1/2 are totally arbitrary.
Let us now consider a typical case, suppose there is an electron in prin-
cipal orbital or shell defined by n = 2. We shall now workout what possible
combinations of quantum numbers this electron may have.
Possible values of azimuthal quantum number l that this electron may have are:
l = 0 and l = 1.
Since there are two possible values of l, there will be two sets of values for
magnetic quantum number m l .
The set corresponding to l = 0 will have only one value m l = 0.
The set corresponding to l = 1, m l may have three values: m l = − 1, 0 and
+ 1.
Now corresponding to each set of values of n, l and m l , spin quantum number
ms may have two values: + 1/2 and − 1/2.
Table 1.9 lists the sets of quantum numbers n, l, m l and m s for principle orbital
of n = 2, such that at least one of these quantum numbers is different. If
orbital − 2 has a single electron then it may have one of the eight different
sets of quantum numbers. Each set defines a sub-orbital or sub-shell within
principle shell-2. Table 1.9 tells that second principle shell (n = 2) has eight
sub-shells. Since quantum numbers associated with each sub-shell are different,
a maximum of eight electrons may be accommodated in second principle shell
(Pauli’s exclusion principle).
Table 1.9 Possible sets of different quantum numbers in orbital (shell) of principal quantum number
n=2
Serial Principal quantum Azimuthal quantum Magnetic quantum Spin
number number n number l number ml quantum
number ms
1 2 0 (s) 0 + 1/2
2 2 0 (s) 0 − 1/2
3 2 1 (p) −1 + 1/2
4 2 1 (p) −1 − 1/2
5 2 1 (p) 0 + 1/2
6 2 1 (p) 0 − 1/2
7 2 1 (p) +1 + 1/2
8 2 1 (p) +1 − 1/2
It is easy to show that the first principal orbital (shell) has only two sub-shells
with set of quantum numbers (n = 1, l = 0, m l = 0, m s = +1/2) and (n =
1, l = 0, m l = 0, m s = −1/2). It is left as an exercise to show that the third
principal shell (n = 3) will have 18 sub-shells and that fourth principal shell 32
sub-shells.
It follows from above that a maximum of two electrons can be accommodated
in I-principal shell, a maximum of 8 electrons in II-principal shell, a maximum
of 18 electrons in III-principle shell, a maximum of 32 electrons in IV-principal
shell and so on.
The principle shells or orbitals are also called electron energy levels and sub-
shells as electron energy states. In n = 1 energy level there are two possible
energy states. Similarly in n = 3, energy level there will be 18 energy states.
1.2.3 Shape and Orientation of Orbitals
Orbital is a three-dimensional space round the nucleus where there is large probability
of finding an electron. There may be several orbitals like, 1s (n = 1, l = 0), 2s (n
= 2, l = 0), 1p (n = 1, l = 1), 3d (n = 3, l = 2) etc.. All these orbitals have
different shapes. Radial probability distribution of electron in hydrogen atom in
1s orbital is shown in Fig. 1.18a. As may be observed in this figure, probability
of finding the electron sharply increase with radial distance, reaches a maximum at
around 0.1 nm from the nucleus and then starts dropping sharply, touching a very low
value at around 0.2 nm and then becomes almost zero with in a small distance. The
probability distribution for 1s orbital is symmetrical in all directions and, therefore, it
appears spherical in 3-D space; the surface boundary diagram of 1s orbital is shown
in Fig. 1.18b where it may be observed that almost 95% chance of finding the electron
is in a spherical volume lying from 0.08 to 0.17 nm from the nucleus. In Fig. 1.18b
the darkness of the colour shade indicates the probability, darker the colour higher
the probability.
Radial probability distribution of electron for 2s orbital (n = 2, l = 1) is shown
in Fig. 1.19a. Note that in this case probability increases from r = 0 and attains a
small maximum value at around r = 0.05 nm and then falls sharply to zero at about r
= 0.1 nm. After touching zero value probability again raises and attains a maximum
value at around r = 0.28 nm and then falls off sharply. In contrast to the case of 1s,
radial probability distribution for 2s orbital shows two maximums, one smaller and
the other larger. In the region between these two maximums probability of finding
electron is zero. This region with probability zero is called the node. Like 1s orbital,
radial probability distribution for 2s is also same in every direction. Boundary surface
diagram of 2s orbital is given in Fig. 1.19b.
Radial probability distribution for 2p orbital, shown in Fig. 1.20, is not symmet-
rical; it has different shapes along the X-, Y- and Z-directions. Though probability
distribution has two bob structure in each direction, the orientation of these bobs
Fig. 1.18 a Electron probability distribution as a function of distance from the nucleus for orbital
1s. b Boundary surface diagram for 1s orbital
Fig. 1.19 a Radial probability distribution for 2s orbital. b Boundary surface diagram for 2s orbital
is different. These different orientations result from different values of magnetic

quantum number m l . It can be shown that probability distributions for orbital-d will
have five different orientations, orbital-f, 7 different orientations and so on.
Fig. 1.20 Three different orientations of electron probability distributions for 2p orbital
1.2.4 Electron Energy Level Diagram
While introducing quantum numbers associated with electron, it was stated that
principal quantum number ‘n’ essentially defines the energy of the electron. Elec-
trons with principal quantum number n = 1 mean electrons that are nearest to the
nucleus; most tightly bound to the nucleus, having largest negative binding energy
and minimum absolute energy. Electrons with n = 2 are not as close to the nucleus
as n = 1 electrons, have negative binding energy but less than that of n = 1 elec-
trons; have absolute energy more than that of electrons of n = 1 orbital. It may thus
be observed that actual (negative) binding energy of electrons decreases, absolute
energy and distance from the nucleus increases as the value of principle quantum
number ‘n’ increases. One may, therefore, draw an energy level diagram for elec-
trons on a vertical scale as shown in Fig. 1.21. It may, however, be said that this
energy diagram is based on quantum mechanical solution of Schrodinger equation
for hydrogen atom.
Fig. 1.21 Energy level

diagram for electron in
hydrogen atom
The fact that a given principal orbital or major shell contains sub-orbitals or sub-
shells brings out the fact that an electron in different sub-shells of a given major shell
will have different energies. It means that the energy of an electron is decided not
only by the principle quantum number ‘n’, but it also depends on the value of the
azimuthal quantum number ‘l’. Actually, an electron put in different sub-shells of a
given major shell will have slightly different energies in different sub-shells.
1.2.5 Electron Configuration of Elements
Though exact solution of Schrodinger equation is possible only for hydrogen atom,
electron energy level scheme obtained for hydrogen atom may be extended, with
some modifications, to obtain the electron configuration for atoms of other elements.
By electronic configuration one means how electrons are distributed in different
orbitals in an atom. From the analysis of hydrogen atom, it is known that principal
orbital with n = 1 may accommodate a maximum of two electrons in sub-state 1s;
principal orbital n = 2, a maximum of 8 electrons (two electrons in sub-state 2s and
6 electrons in sub-state 2p) and so on. The problem in case of atoms other than that
of hydrogen is to find out the sequence of sub-orbitals with increasing energy.
1.2.6 Aufbau or Building-Up Principle
The underlying principal of electronic configuration is that electrons in different

orbitals of an atom are distributed in such a way that the atom has minimum energy.
This means that electrons in an atom are filled in different orbitals in order of their
increasing (absolute) energy. This means that first electrons are filled in orbitals that
have minimum absolute (or maximum negative binding energy) energy. Once orbital
with minimum absolute energy is filled, additional electrons will go to the orbital or
state that has next higher value of absolute energy and so on. As already mentioned,
the energy of orbitals increases with the value of principal quantum number ‘n’ in
case of Hydrogen atom, but in case of other atoms that have more than one electron
energy sequencing of orbitals is found to depend not on ‘n’ but on the sum (n +
l). Following rules may be used to determine the energies of different orbitals and
sub-orbitals.
RULE-1: An orbital with a lower value of (n + l) has lower energy. For example
let us consider two orbitals 4s and 3d . The value of (n + l) for 4 s is (4 + 0 = ) 4
and for 3d (n + l) = (3 + 2) = 5. Thus orbital 4s will have lower energy and will be
filled before orbital 3d .
RULE-2: If the value of (n + l) for two orbitals is same then orbital with lower
value of n will have lower energy and will be filled before the other. For example
consider the orbitals 3d and 4p the value of (n + l) for orbitals 3d = (3 + 2) = 5;
and for orbital 4p (n + l) = (4 + 1) = 5. So both these orbitals have same value of

(n + l), hence according to RULE-2, orbital with lower value of n, that is 3d will
have lower energy than that of orbital 4p.
RULE-3: Hund’s rule: This rule concerns the distribution of electrons in a set
of orbitals of same energy, i.e. in sub-shells or sub-orbitals. According to this
rule if a number of orbitals of the same sub-shell are available then electrons
distribute in such a way that each orbital is first singly occupied with same
spin. For example consider the distribution of electrons in carbon atom. A carbon
atom has six electrons, now first two electrons will be accommodated in the lowest
energy orbital 1s; one electron each in sub-orbitals 1s+1/2 and 1s−1/2 . The remaining
four electrons will go to orbitals 2s and 2p. Out of these four electrons two will
be accommodated in states 2s+1/2 and 2s−1/2 . The remaining two electrons will now
go to 2p orbital. Now 2p orbital may have six sub-orbitals 2px+1/2 , 2px−1/2 , 2py+1/2 ,
2py−1/2 , 2pz+1/2 and 2pz−1/2 . Hund’s Rule says that the two electrons that will be
accommodated in 2p orbital will not go to 2px , but will distribute such that one
electron will go in state 2px+1/2 and the other in 2py+1/2 . Since there were only two
electrons, state 2px−1/2, 2py−1/2 , 2pz+1/2 , 2pz−1/2 will remain unfilled.
According to the rules mentioned above the sequence of orbitals with increasing
energy comes out to be:
1s < 2s < 2p < 3s < 3p < 4s < 3d < 4p < 5s < 4d < 5p < 6s.
Energy sequence mentioned above is shown in Fig. 1.22, where energy of orbitals
in increasing order is shown.
An easy way of remembering the energy sequence of orbitals is shown in Fig. 1.23.
So, the energy sequence of orbitals may be listed as:
1s, 2s, 2p, 3s, 3p, 4s, 3d, 4p, 5s, 4d, 5p, 6s, 4f, 5d, 6p, 7s, 5f ….
Fig. 1.22 Energy level

scheme of orbitals
Fig. 1.23 Sequencing of

orbital energies
1.2.7 Representing Electron Configuration
With the help of the rules discussed above it is possible to write the electron config-
uration for atom of any element. There are three ways of representing electron
configurations.
(a) Orbital notation method In this method the orbitals that have electrons are
written in order of increasing energy and the number of electrons in each
orbital are given as a superscript to the orbital. For example, nitrogen atom has
seven electrons and its electronic configuration may be written as: 1s2 2s2 2p3 .
Figure 1.24 explains the meanings of each character of the notation.
(b) Orbital diagram method In this method orbitals having electrons are repre-
sented by boxes and are written in the order of increasing energies. Electrons in
each orbital are represented by arrows, direction of arrows indicating the direc-
tion of electron spins. For example, the electron configuration of some elements
atoms are given in the last column of Table 1.10.
Fig. 1.24 Meaning of each

character of the notation
Table 1.10 Electronic configurations for some elements in different notations
(c) Short-hand form In this method the last completely filled orbital or shell is
represented in terms of a noble gas. For example, the electron configuration of
lithium in this notation may be written as [He] 2s1 . Electron configurations for
some elements in different notations are given in Table 1.10.
1.2.8 Valence Shell
Shell or orbital of highest energy (largest value of n) that has some electrons is
called the valence shell or valence orbital, and the electrons it contains are called
valence electrons. For example the valence shell for Calcium is 4s2 with two valence
electrons; the valence shell for Argon is 3s2 3p6 orbital with (2 + 6 =) 8 valence
electrons; Aluminium has 3s2 3p1 shell as valence shell and it has (2 + 1 =) 3 valence
electrons. It is important to remember that all electrons in different sub-shells of the
highest ‘n’ value shell (that has some electrons) are counted as valence electrons. The
number of valence electrons in case of noble gases is eight. Therefore, it is concluded
that eight electrons in valence shell of any atom make it very stable and chemically
inert.
Valence shell and valence electrons are important because most of the chemical
and some physical properties of the atom are decided by the valence shell and valence
electrons. It is valence electrons that take part in chemical reactions and decides the
type of bonding with other atoms to make molecules.
Each orbital or shell/sub-shell can accommodate a fixed maximum number of

electrons; for example, s-shell can accommodate a maximum of 2 electrons, p-shell
a maximum of 6 electrons, d-shell a maximum of 10 electrons and f-shell a maximum
of 14 electrons. When the valence shell of an atom is filled with maximum number
of electrons it can hold, the valence shell is said to be completely filled and the atoms
with completely filled valence shell are found to be chemically very stable and do not
show chemical activity. Valence shells of inert gases like argon, neon, and helium,
etc. are completely filled with eight electrons.
In periodic table (see Fig. 1.1) elements are arranged such that all elements falling
in one group (vertical column) have similar electron configurations of their valance
shells. They have same number of valence electrons; for example configurations of
members of group-II are given as: Be = [He] 2s2 , Mg = [Ne] 3s2 , Ca = [Ar] 4s2 , Sr
= [Kr] 5s2 , Ba = [Xe] 6s2 , and Ra = [Rn] 7s2 . Each member of group-II has only
two valence electrons in s-orbital. This is the reason that all members of a group of
periodic table show similar chemical properties.
The periodic table may also be taken as a guide to the order in which orbitals are
filled. Figure 1.25 shows the classification of groups of elements in the periodic table
according to the type of outer sub-shell being filled with electrons.
In Fig. 1.25 of the periodic table different elements may be identified through their
horizontal group number given in red colour and the vertical row number written in
blue. For example, the element specified as 75 is 43Tc; 46 (coloured dark brown)
is 72Hf and element specified as 125 (shown by light brown colour) is 48Cd. The
important characteristic of this figure of periodic table is that all elements falling
in a group (vertical column) have same general formula for the valence orbital.
For example all members of group 1A the general formula for valence orbital is:
ns1 where n is the principle quantum number. All members of group 2A have the
common representation for the valence orbital as: ns2 ; for all members of group 3A,
the general representation for valence orbital is: ns2 np1 ; … for members of the group
6A: ns2 np4 and so on. The outer electron configuration for element 72Hf is: 6s2 5d2 ;
Fig. 1.25 Groups of

elements in periodic table
according to the type of
sub-shells being filled with
electrons
similarly, the valence electron configuration for element 48Cd is: 5s2 4d10 and for
the element 43Tc the outer electron configuration is: 5s2 4d5 .
1.2.9 Some Anomalous Electron Configurations
The electron configuration rules stated above holds good in most cases but there
are four outstanding exceptions where these rules fail to give the correct electron
configuration. These four cases are of:
(a) Chromium Cr; According to the rules the electron configuration of Cr should
be [Ar] 4s2 3d4 but actually it is [Ar] 4s1 3d5
(b) Similarly, for copper Cu, according to rules the electronic configuration should
be[Ar] 4s2 3d9 but its actual configuration is: [Ar] 4s1 3d10
(c) Silver (Ag) according to rules should have electron configuration of [Kr] 5s2
4d9 but actual electron configuration for silver is: [Kr] 5s1 4d10
(d) Also in case of Gold (Au) according to rules the electronic configuration should
be [Xe] 6s2 5d9 but the actual configuration is: [Xe] 6s1 5d10
As may be observed in all the above cases, an enhanced stability is acquired by
half or fully filled sub-shells.
Only valence electrons take part in chemical reactions and in forming molecules,
etc. The inner electrons are generally well protected and mostly do not take part
in combination processes. An American chemist, G. N. Lewis introduced a simple
notation to represent valence electrons in an atom. These notations are called Lewis
symbols. Lewis symbols for elements of second period of periodic table may be
given as (Fig. 1.26).
The dots around the chemical symbol of an element give the number of valence
electrons in the atom of the element. The valency of the element is equal to the
number of dots around it or 8 minus the number of dots around.
Electron configurations discussed above apply to the ground states of atoms.
However, when an atom is excited by giving some energy by an external source, say
by heating, etc., few electrons from its valence orbital shift to the next higher orbital.
Therefore, the electron configuration of an excited atom is different from the ground
state configuration. Similarly, electron configuration of an atomic ion is different
from that of the parent atom.
SAQ: What is the difference between a classical electron orbit and quantum
mechanical orbital?
Fig. 1.26 Lewis symbols for second group elements

1.3 Bonds Between Atoms and Ions 43
SAQ: Does the energy of all electrons in a given principal orbital exactly same?
SAQ: Which quantum number includes Pauli’s exclusion principle in quantum
description of electron’s motion in an atom?
SAQ: Write electron configurations of a singly ionised Sodium ion and a doubly
ionised lithium ion.
1.3 Bonds Between Atoms and Ions
When two atoms are brought near to each other two types of Coulomb forces come
into play; the repulsive forces between the positively charge nuclei of the two atoms
and between their electron clouds and attractive forces between electron cloud of
one atom and the nucleus of the other atom. The magnitude of both types of forces
increases with the decrease of relative separation r. Figure 1.27a shows the variation
of the attractive, repulsive and net forces between two atoms as a function of their
mutual separation r. The intra-atomic separation r 0 corresponds to the mutual separa-
tion where attractive and repulsive forces cancel each other and the two atoms are in
a state of equilibrium. It is well known that any force F may be converted into poten-
tial energy V (or force F may be derived from potential V ) using the mathematical
operation given by expression;
( )
∂F ∂F ∂F
V = −gradF = − + +
∂x ∂y ∂z
Net potential energy between the two atoms as a function of intra-atomic sepa-
ration r, obtained by using the above expression, is shown in Fig. 1.27b. It may
be observed in this figure that net attractive force gives rise to the negative poten-
tial energy V 0 which is responsible for the binding of the two atoms. The negative
potential energy has its maximum (negative) value at a separation r0 , the equilibrium
distance. Two atoms develop a bond only when they are at a relative separation of r 0 ,
at larger separation they do not bind with each other as shown in the figure. The nega-
tive binding energy decreases on both sides of r 0 , and, therefore, the two atoms are
held fixed at a separation of r 0 . Further, larger the magnitude of − V 0 , more tightly
the two atoms are bound with each other. Atoms bound with each other makes a
molecule. As a matter of fact the magnitude of V 0 decides many properties, physical
and chemical, of the pair of two atoms or the molecule. For example, a larger value
of V 0 corresponds to a higher melting point.
1.3.1 Electronegativity
Chemical reactivity of an atom is decided by the valence electrons in the outer most
orbital of the atom. If the valence shell is completely filled, like that of inert gases,
Fig. 1.27 a Attractive,

repulsive and net force
between two atoms as
function of intra-atomic
distance r. b Net potential
energy between two atoms as
a function of r
the atom has almost no chemical reactivity. However, in case the valence shell is only
partially filled, the atom shows chemical reactivity. Chemical reactivity of atoms may
be measured in terms of parameters called electronegativity or electropositivity.
Figure 1.28 shows the periodic table of elements. The electronegativity of elements
in periodic table starts from almost the middle of the periodic table and increases
towards right. Electronegativity decreases towards left, and atoms with lower value
of electronegativity are said to have electropositivity. In periodic table highly elec-
tronegative halogens and highly electropositive alkali metals are separated by the
noble gases. Atoms of elements of group-1 and group-2 of the periodic table (enclosed
in box A) have partially filled valence shell; they have only one or two electrons in
their valence shell. These atoms having almost empty valence shell, with lower elec-
tronegativity, have the tendency to give up their electrons to other atoms of higher
electronegativity when they come in contact with them. Since they have the tendency
to give their electrons and by doing so they acquire net positive charge, these atoms
are said to have electropositivity. For example, if we consider element potassium
(K) its electron configuration is: 1s2 2s2 2p6 3s2 3p6 4s1 or [Ar]4s1 . The valence shell
of potassium 4s1 contains only on electron. Potassium has the tendency to give up a
electron and becomes positive ion;
K − 1 electron = K+ positive ion (Cation).

Fig. 1.28 Periodic table showing electronegative and electropositive elements
On the other hand atoms of elements on right side of the periodic table like,
chlorine (Cl) has many electrons in its valence shell but is not completely filled.
Electron configuration of chlorine is: [Ne] 3s2 3p5 . The valence shell of chlorine
(3s2 3p5 ) has seven electrons, one electron short of the maximum number that (sp)-
orbital can accommodate. Chlorine is highly electronegative; it has the tendency of
taking an electron and becoming a negative ion (anion);
Cl + 1 electron = Cl− Negative ion (anion).
Elements contained in box A in Fig. 1.28 are typical of metallic character (elec-
tropositive), and those contained in box B have characteristics which are intermediate
between metals and non-metals, possessing different degrees of electronegativity.
1.3.2 The Octet Rule
Though several attempts were made to explain the bond formation between atoms
on the basis of electron structure, it were Kossel and Lewis, who independently gave
a satisfactory explanation in 1916. They studied the electron structure of noble gases
and found that all of them have eight electrons in their valence shell. Based on this
observation, Kossel and Lewis developed a theory for combination between atoms
called ‘electronic theory of chemical bonding’. According to this theory atoms
can combine either by transfer of valence electron from one atom to the other or
by shearing of valence electrons in order to have an octet (eight) of electrons in
their valence shells. This is called octet rule. The octet rule though useful but is not
universal. There are some limitations of the octet rule. Octet rule essentially applies
to atoms of the second group of periodic table.
1.3.3 Classification of Bonding
It is clear from above that characteristics of valence electrons and the net force of
attraction between two atoms create attractive bonds between atoms. Depending on
their strength and other characteristics bonds between atoms may be divided into
two classes: (a) primary bonds between atoms and (b) secondary bonds between
atoms and molecules.
Primary bonds may further be divided into three types: (i) ionic bond, (ii) covalent
bond and (iii) metallic bond while secondary bonds are into two types: (i) van der
Waals bond and (ii) hydrogen bond.
(A) Primary atomic bonds Primary atomic bonds are characterised by large inter-
atomic forces. Primary bonds involve valence electrons of interacting atoms
and arise from the tendency of atoms to acquire stable electron structure of
completed valence shell. They may be nondirectional or localised (directional)
and may be produced by electron transfer, electron shearing or delocalisation
of electrons
(i) Ionic or electrovalent bond These bonds are formed when two atoms
of very different values of electronegativity combine together to form a
molecule. Ionic bonds are formed typically between highly electropositive
(metallic) and electronegative (non-metallic) elements. For example when
an atom of sodium with very low value of electronegativity (or high value
of electropositivity) combines with an atom of chlorine which has very
high value of electronegativity, an ionic bond gets established between
the two atoms.
Fig. 1.29 Formation of ionic bond in NaCl molecule
As shown in Fig. 1.29, sodium (Na) has only one electron in its valence
shell 3s. Since sodium is electropositive it has the tendency of giving
away this electron, chlorine on the other hand is highly electronegative,
has 7 electrons in valence shell 3s2 3p5 and has the tendency of acquiring
electrons. As a result, when an atom of sodium comes sufficiently close
to the chlorine atom to be with in its Coulomb field, it gives its only
valence electron to the chlorine atom. With this transfer of electron sodium
atom becomes a positive ion and on receiving an extra electron from
sodium, chlorine atom becomes a negative ion. A bond gets established
between positive sodium ion and negative chlorine ion due to coulomb
attraction between them. Thus ionic bonds are formed by the transfer of
valence electron(s) from the lower electron negativity atom to the higher
electronegative atom.
Only valence electrons take part in bonding while the inner shell electrons
and the nucleus of the atom do not take part in chemical bonding, being
well shielded by valence electrons and large force of attraction between
the nucleus and inner shell electrons. Therefore, in pictorial representa-
tions of bonds, the nucleus and inner shell electrons are often represented
by the core. Core of the atom has a positive charge equal in magnitude
of the number of electrons in the valence shell. Figure 1.29 explains the
ionic bond formation in case of NaCl. As shown in this figure, on forming
the bond the size of Na+ ion shrinks (as compared to the size of Na atom)
while the size of negative Cl− ion also shrinks but not so much as that
of sodium ion. The reason for this reduction of size in case of sodium
ion is the fact that on losing the only valence electron the valence shell
disappears. The size of sodium ion then reduces to the size of its core. The
size of negative chlorine ion also decreases because with the increase of
Fig. 1.30 a Ionic bonds are generally non directional. b Structure of an ionic solid
negative charge in the valence shell, attractive force by its core increases
which results in a valence orbital of reduced size.
Ionic bonds do not have any preferred direction; it is because of the
fact that both the positive ion and the negative ion attract each other by
forces of equal magnitudes, as shown in Fig. 1.30a. Ionic solids have large
lattice energies ranging from 600 to 3000 kJ/mol and have high melting
temperatures. Melting point for NaCl, for example, is 801 °C. The binding
energy for NaCl is ≈ −7.42 × 10−19 J = 4.63 eV.
Ionic solids have a regular arrangement of alternate positive and negative
ions in three dimensions, as shown in Fig. 1.30b. Ionic solids are mostly
ceramics; they are bad conductors of heat and electricity and often brittle.
(ii) Covalent bonds
The term covalent bond was coined by the American Chemist Irving Lang-
muir in 1919. Covalent bonds are formed when two atoms of either same
or nearly same electronegativity join together. This typically happens
in non-metals. In case of covalent bonding, transfer of electrons from one
atom to the other does not take place. Instead, the two interacting atoms
shear electrons to complete the octet (eight electrons each) or duplet in
their valence shells.
Large number of elements in periodic table have either s-shell or the
combination of sp orbitals as their valence shell. The maximum number
of electrons that may be accommodated in s-shell is 2, and for sp shell
8 (2 + 6), therefore, atoms try to complete electron duplet, if s-shell is
valence shell or octet if sp shell is valence shell.
Fig. 1.31 a Covalent bonding in Cl2 molecule. b Lewis dot structure for Cl2 molecule
Shearing essentially means that some electrons of individual atoms get

attached to both the atoms. Generally, s and sp shell electrons shear to
attain noble gas electron configuration. Covalent bonds are formed when
atoms of F, O, N, Cl, H, C, Si, Ge, etc. form molecules. Covalent bonds
are also found in compounds like GaAs, CH4 (methane), C2 H6 (ethane),
etc.
Figure 1.31a shows how two atoms of chlorine with seven electrons each
in their valence shells share one electron to complete octets resulting in
stable Cl2 molecule. Lewis dot structure of Cl2 molecule is shown in
Fig. 1.31b.
In the above-mentioned example of Cl2 molecule, two atoms of chlorine
share only one electron creating a single covalent bond. However an atom
may have more than one covalent bonds. For example, in one molecule
of methane (Ch4 ) carbon atom that has only four valence electrons share
an electron each with four hydrogen atoms, creating four covalent bonds.
This is shown in Fig. 1.32. The number of covalent bonds that an atom
may form is given by (8 − N) where N is the number of valence electron
in the atom. That is why a carbon atom with 4-valence electrons forms
(8 – 4 =) 4 covalent bonds and Chlorine with seven valence electrons
(8 – 7 =)1 covalent bond.
In C2 H4 molecule the two carbon atoms share two electrons each to
complete electron octet for both carbon atoms. Thus the two carbon atoms
get coupled through a covalent bond doublet. The remaining two electrons
of each carbon atom are sheared by single electrons of two hydrogen atoms
per carbon atom, completing the duplet of hydrogen.
Fig. 1.32 a Pictorial representation of four covalent bonds in methane (Ch4 ) molecule. b Lewis
dot structure for methane molecule
Lewis dot structure of C2 H4 molecule is given in Fig. 1.33a. Nitrogen

atom with 5 electrons in s2 p5 valence shell is unstable, but it makes a
stable N2 molecule with three covalent bonds as shown in Fig. 1.33b.
Bond Parameters
(a) Bond length Bond length is defined as the equilibrium distance
between the nuclei of two bonded atoms in a molecule. Each atom
of the molecule contributes to the bond length. Figure 1.34a shows
two atoms A and B that are bonded by covalent bound. Here R is the
covalent bond length and the contributions of the two atoms to the
bond length, r a and r b are the bound radii respectively of the two
atoms. The covalent bond lengths for O–H, C–H, C=C molecules
are respectively 96, 107, 120 (× 10–12 m).
Fig. 1.33 a Lewis dot structure of C2 H4 molecule with double covalent bonds in carbon atoms.
b Lewis dot structure for N2 molecule with triple covalent bonds in nitrogen atoms
Fig. 1.34 a Covalent radii of the two atoms are r a and r b , while R the separation between the two
nuclei is the bond length. b Bond angle in H2 O molecule
(b) Bond angle It is defined as the angle between the orbitals containing
bounding electron pair around the central atom in a molecule or
complex ion. Bond angle is expressed in degrees, and it gives some
idea about the distribution of orbitals around the central molecule,
which means the shape of the molecule. For example, in case of
water molecule the bond angle is 104.50, as shown in Fig. 1.34b.
(c) Band order In Lewis description of covalent bond, bond order is the
number of covalent bond in the molecule, for example in H2 bond
order is 1, in O2 bond order is 2 and in N2 the bond order is 3.
(d) Polarity of bond When covalent bond is formed in two similar atoms
like H2 , O2 , Cl2 , etc., the shared electrons are equally attracted by
the two atoms and, therefore, the electron pair is situated exactly
between the two identical atoms or nuclei. Such a bond is called
a non-polar covalent bond. However, when two dissimilar atoms
are coupled through a covalent bond, like HF molecule, the sheared
pair of electrons shifts more towards the fluorine because of the
larger force of attraction by it than that by hydrogen nucleus. The
bond in this case is a polar covalent bond. Shifting of the paired
electrons from the centre gives rise to the formation of an electric
dipole. The molecule that has polar covalent bond behaves as a tiny
electric dipole, often represented by the symbol μ. Dipole strength
of such atoms is measured in a unit called Debye denoted by D.
Further 1 D = 3.33564 × 10–30 C m. Dipole moment is a vector
quantity, and the direction of the vector is indicated by direction
of shift of the shared electron pair from the central position, i.e.
in case of HF molecule from H atom towards F atom. Because of
the associated dipole moment polar covalent bonds are said to have
specific directions.
(e) Bonding energy The bonding energies of covalent bonds may be
very different; it may be very high, for example in case of diamond
which is the hardest material having a melting point of > 3550 °C.
On the other hand bonding energy may be very low as in the case of
bismuth which has a low melting point of around 270 °C.
No material has 100% ionic bonds or 100% covalent bonds; materials
that have predominantly ionic bonds also have a small percentage of
covalent bonds and vice versa.
(iii) Metallic bond
Metallic bonds are found in metals and their alloys. Such bonding occurs
when atoms of low electronegativity join together. Since low electroneg-
ativity atoms have the tendency to lose their valence electrons; the inter-
acting atoms lose all their valence electrons which form a cloud of delo-
calised electrons. No electron of the cloud is essentially attached with
any particular atom rather all electrons are attached with positive cores
of all atoms. Metallic bonding may be looked as an extreme case of
covalent bonding; in covalent bonding nearby atoms shear their valence
electrons but in metallic bonding all atoms shear their valence electrons.
Electrons that are not bound to any particular atom are called delocalised
electrons. Cloud of delocalised electrons works as glue to bind positively
charged cores of atoms, which in absence of delocalised electron cloud
will repel each other and break the crystalline structure of the metal. Alter-
nately, a metallic crystal that has N number of atoms may be looked as
an N-atomic molecule and the cloud of delocalised electrons as electrons
moving in large number of different molecular orbitals. Since the number
of molecular orbitals of a molecule having N atoms will be very large,
criss-crossing each other, the delocalised electrons in these molecular
orbitals appears as an electron cloud.
Looking from the point of quantum mechanics, each electron of the
electron cloud has a discrete set of closely spaced energy levels. Since
the number of delocalised electrons in the cloud is very large, the total
number of electron energy levels becomes very large, almost a continuum
of levels. Only the low lying energy levels of the continuum are filled
with electrons, but large number of electron levels is empty. If energy is
supplied to these delocalised electrons by some external source, say by
shining light on the metallic crystal, electrons absorb the incident light
photons of all frequencies and go to their respective higher energy levels
which were empty. Since the mean life of these excited states is very short
(< 10−9 s), the excited electrons revert back to their lower energy states
re-emitting photons of almost the same energy as they have absorbed.

This is why freshly cut metallic surfaces shine when put in light.
Metallic bond energies may vary by large factor, for example, the bond
energy in case of Tungsten is 850 kJ/mol and its melting point is 3410 °C.
While bond energy in case of mercury is only 68 kJ/mol, with melting
point of − 39 °C.
Since the force responsible for binding atomic cores in a metallic crystal
is provided by the cloud of delocalised electrons that does not have a fixed
shape, it is relatively easy to bend metals. Further, since ionic cores of
metallic atoms may slide over each other, metals are ductile. Some details
of these properties of metals have already been discussed earlier. Because
of close packing of atoms, most of the metallic crystals are solid (except
mercury), have high density, high melting point, etc. The electron cloud
of nearly free delocalised electrons makes metals a good conductor of
both electricity and heat.
(B) Secondary bonds
Secondary bonds, as their name suggests, are relatively weaker than primary
bonds, may have strength of the order of 10 kJ/mol, but they are very important.
Secondary bonds, also called van der Waals bonds, are present in most of the
systems but are overlooked because of their low strength. Secondary bonds
originate because of electric dipole nature of some molecules/atoms.
Electric dipole Fig. 1.35a shows an electric dipole; two equal and opposite
charges separated by a small distance, and the vector dipole moment represented
by Greek letter μ = Q . x. Here Q is the magnitude of charge (in Coulomb)
and x (in metre) is the distance between the two charges. Dipole moment is a
vector along the direction from negative charge to the positive charge.
(i) Fluctuating dipole bond Electric dipoles are often formed in micro-
scopic systems like atoms and molecules if centres of positive and nega-
tive charge distributions in them do not coincide. A charge distribution
if looked from a distance appears as if the total charge contained in the
distribution is concentrated at a particular point, this point is called the
charge centre. This is likely to happen in a molecule/big atom, the electron
clouds of which may extend up to large distances from their nuclei and
may have non-spheroid shape. As a result, the centre of positive charge
gets slightly displaced with respect to the centre of negatively charged
electron cloud. Such molecules/atoms are called polar molecules/atoms.
Polar molecules or atoms have the tendency to align themselves such
that their electric dipole moment vectors point in opposite directions.
Oppositely directed dipoles attract each other creating a secondary (or
van der Waals) bond, as shown in Fig. 1.35b. Molecular dipoles formed
as a result of displaced charge centres are often referred as fluctuating
dipoles. Typical examples are atoms of noble gas elements. Noble gas
Fig. 1.35 a An electric dipole of dipole moment μ. b Attraction between dipole molecules/atoms
give rise to secondary bonding
atoms have completely filled valence shell (s2 p6 ) and therefore, cannot
form primary bonds. However, when two atoms of a noble gas come close
to each other, they induce electric dipoles in each and these dipoles align
to form a fluctuating dipole secondary bond, as shown in Fig. 1.36a.
(ii) Permanent dipole bond However, there are molecules that are permanent
dipoles, like that of NaCl, which have ionic bonding between Na+ and Cl−
ions. In such molecules close packing and alignment results in formation
of permanent dipole secondary bonds (Fig. 1.36b).
(iii) Hydrogen (secondary) bond Hydrogen bond is also a secondary dipole
bond but it is much stronger compared to other secondary bonds.
Hydrogen bond energy may be as large as 50 kJ/mol. This bond is found
Fig. 1.36 a Fluctuating dipole secondary bond b permanent dipole secondary bond
in molecules where hydrogen is covalently bound to elements of high

electronegativity like, fluorine, oxygen or nitrogen.
When hydrogen is bonded with covalent bond to strongly electronegative
element ‘x’, the electron pair sheared between hydrogen and atom ‘x’, moves
faraway from hydrogen atom. Since there is displacement of electrons towards
‘x’, the hydrogen atom acquires a small positive charge (q+) and the other
atom ‘x’ a negative charge (q−). This results in molecule becoming polar and
behaving like a dipole. Two polar molecules then make a secondary (hydrogen)
bond. Covalent and hydrogen bonds between molecules of HF are shown in
Fig. 1.37. It is because of the relatively strong hydrogen bonds between water
molecules that the boiling point of water is so high as compared to other
materials of same molecular weight.
Figure 1.38 shows how hydrogen bonds bind large number of water molecules
together. Each water molecule (H2 O) has two hydrogen atoms that are coupled to
oxygen atom with covalent bonds. Each water molecule behaves as a dipole since the
pairs of electrons in covalent bonds are slightly shifted towards the oxygen nucleus
so that hydrogen atoms develop positive charge and oxygen atom negative charge.
The dipole water molecules join together by hydrogen bonds. Often in pictorial
representation, covalent bonds are shown by solid line and hydrogen bond by dashed
line, as shown in part (b) of the figure.
Properties of six types of bonds are summarised in Table 1.11.
SAQ: Primary cause of bound formation between atoms is?
SAQ: What works as a glue between positive cores in crystalline metals?
SAQ: Give one similarity between covalent and metallic bonding.
SAQ: What is the difference between fluctuating and permanent dipole bonds?
SAQ: Hydrogen bond is a fluctuating dipole or permanent dipole bond?
Short Answer Questions
SA 1.1 Two metals A and B respectively, have face centred (fcc) and body centred
(bcc) crystal structures. Which of the two metals will be more ductile and
why? (II).
Fig. 1.37 Hydrogen bond

between two molecules of
HF
Fig. 1.38 Hydrogen bonds between water molecules
Table 1.11 Summary of bond properties

Bond class Special characteristics Typical bond Examples
energies
Ionic Formed by electron transfer between two 300–6000 kJ/ NaCl, CaCl, LiF
atoms of very different electronegativity, mol or
non-directional, very strong 5–10 eV/atom
Covalent Formed between atoms of either same or 300–800 kJ/ Ge, Si, CH4 ,
nearly equal electronegativity, shearing of mol or 3–8 eV/ Diamond
electrons, may be directional if two atoms are atom
different
Metallic Delocalisation of valence electrons, 100–1000 kJ/ All metals and
delocalised electron cloud works as glue to mol or their alloys
hold positive cores, good conductors of heat 0.5–1.6 eV/
and electricity, non-directional atom
Dipole Formed by alignment of either fluctuating or Around 10 kJ/ Noble gas
bond permanent dipoles of molecules/atoms, mol or molecules
non-directional 0.05–0.2 eV/
atom
Hydrogen Relatively stronger secondary bonds Around 50 J/ H2 O, HF, HN
bond mol or
0.25–0.6 eV/
atom
SA 1.2 What are delocalised electrons? Discuss their role in determining the
resistivity of a metal.
SA 1.3 Why the resistivity of a metal increases with temperature?
SA 1.4 Why does a piece of metal become more malleable on heating and tough,
brittle on cooling and hammering?
SA 1.5 How one may define a Ceramic? Are both Diamond and Graphite Ceramic,
if yes, why?
SA 1.6 How can one differentiate between Metals, Ceramics and Polymers on
the basis of bonding in them? Which property of Ceramics is very much
affected by the relative strength of bond types in it?
SA 1.7 What are thermosetting and thermoplastic polymers? Give one example
of each.
SA 1.8 Discuss the process of crack/fracture repair in ceramic fiber reinforced
CMCs.
SA 1.9 What are light metal reinforced (MMC) and where are these used?
SA 1.10 Which Composite you will use for fabricating light but strong car body.
Give reasons for your answer.
SA 1.11 State and explain the rules that govern the distribution of electrons in an
atom.
SA1.12 Magnesium nucleus has 12 protons, write the electronic configuration of
magnesium atom in three notations.
SA 1.13 What will be the multiplicity of f-orbital if electron spin is neglected.
SA 1.14 What is the special feature of elements in a group of periodic table and
how does it effects the chemical behaviour of elements, explain with an
example.
SA 1.15 Electron configuration of 21 Sc is [Ar] 3d1 4s2 . What will be the configura-
tion of 22 Ti?
SA 1.16 How electro negativities of interacting atoms decide the type of bond
between them?
SA 1.17 How does the melting point of a material related to the strength of the bond
in its molecules?
SA 1.18 The HF molecule has a permanent dipole moment or not? Explain your
answer.
SA 1.19 What are delocalised electrons? How do they play role in creating atomic
bonding?
SA 1.20 What type of bonds do you expect between long molecular chains in
polymers?
Multiple Choice Questions

Note: Some of the multiple choice questions may have more than one correct
alternative; all correct alternatives must be ticked for complete answer in such cases.
MC 1.1 Silicon carbide (SiC) and Molybdenum silicate (MoSi2 ) are;
(a) Metals (b) Ceramics (c) used as heat shields (d) used for making
electrodes
ANS: (b, d)
MC 1.2 Which of the following is Elastomers?

(a) Polysulphide (b) Terylene (c) Neoprene (d) Sealing wax
ANS: (c)
MC 1.3 Pickup the thermoplastic polymer (s) from the following
(a) Nylon (b) Terylene (c) Neoprene (d) Sealing wax
MC 1.4 Ceramic fibers used as reinforcement in (CMC) are coated with;
(a) Polyethylene (b) Silicon carbide (c) Boron nitride (d) pyrolytic
carbon
ANS: (c), (d)
MC 1.5 Composite used for making brake linings is;
(a) Light metal reinforced (MMC) (b) Carbon fiber reinforced (PMC)
(c) Carbon fiber reinforced (CMC) (d) Carbide reinforced (MMC)
ANS: (c)
MC 1.6 Which of the following is (are) used for making electric heating elements
and electrodes.
(a) Ceramic ZrO2 (b) Ceramic MoSi2 (c) Polymer neoprene (d) Polymer
Buna-N
ANS: (a), (b)
MC 1.7 Bakelite is
(a) Metal matrix composite (b) Thermosetting polymer (c) Thermo-
plastic Polymer (d) Polymer matrix composite
ANS: (b)
MC 1.8 Which of the following may be used for moderation of fast neutrons in
reactors?
(a) Polymer Bakelite (b) Ceramic BeO (c) polymer neoprene (d)
Ceramic ZrO2
ANS: (b), (d)
MC 1.9 Which of the following is/are not ceramic?
(a) SiC (b) ZrO2 (c) MoSi2 (d) Sealing wax
ANS: (d)
12
MC 1.10 Electron configuration of 6 C is 1s2 2s2 2p2 , the electron configuration
of 14
6 C will be;
(a) [He]2s2 2p2 (b) [He]2s2 2p3 (c) 1s2 2s2 2p2 (d) 1s2 2s2 2p1
ANS: (a), (c)
MC 1.11 Azimuthal quantum number and maximum number of electrons that

may be accommodated in f-orbital are respectively;
(a) 2 and 7 (b) 3 and 7 (c) 2 and 14 (d) 3 and 14
ANS: (d)
MC 1.12 Which of the following represent same electron configurations;
(a) [He]2s2 2p2 (b) [He]2s2 2p3 (c) 1s2 2s2 2p2 (d) 1s2 2s2 2p1
ANS: (a), (c)
MC 1.13 The electron configuration of deuteron atom is;
(a) 1s2 (b) 2s2 (c) 1s1 (d) 1s2
ANS: (c)
MC 1.14 Noble gas atoms have completed valence shells; they are bound by.
(a) Ionic bonds (b) covalent bonds (c) fluctuating dipole bonds (d)
permanent dipole bonds
ANS: (c)
MC 1.15 Diamond, the strongest material has.
(a) Ionic bonds (b) covalent bond (c) hydrogen bonds (d) fluctuating
dipole bonds
ANS: (b)
MC (1.16) Covalent bonds are formed when atoms of
(a) very different electronegativity combine (b) nearly same elec-
tronegativity combine (c) carbon combine (d) hydrogen and fluorine
combine
ANS: (b), (c)
MC 1.17 If bond energy of a material is 4.0 eV/atom, the type of bond and its
melting point are respectively;
(a) ionic, high (b) fluctuating dipole, high (c) hydrogen, high (d)
covalent, low
ANS: (a)
Long Answer Questions
LA 1.1 Summarise important properties of metals that differentiate them from other
materials. Explain how delocalised electrons are formed in metals and the
role they play in deciding thermal and electrical conductivities of metals.
LA 1.2 What are composites and how are they classified? Discuss important
properties and applications of fiber reinforced Ceramic matrix composites.
LA 1.3 How can one define a ceramic? Name two ceramics that have electrical prop-
erties opposite to each other and two that have similar electrical properties.
Bring out differences between metals, polymers, ceramics and composites.
LA 1.4 What are polymers? What types of bond are usually found in polymers?
What is meant by the functionality of a monomer and how does it affect the
structure of polymer?
LA 1.5 Discuss quantum mechanical model of electron configuration in atoms and
explain physical significance of different quantum numbers giving their
possible values. What are the rules for filling electrons in different orbitals?
Explain by giving some example.
LA 1.6 What is meant by electronegativity? How does electronegativity decide
nature of bonds between two atoms? Describe various types of primary
bonds giving examples for each type.
LA 1.7 What are secondary bonds and how do they differ from primary bonds?
Give details of fluctuating and permanent dipole secondary bonds bringing
out the points of difference between them. What kind of secondary bonds
are found in water molecules? What are their special characteristics?
Chapter 2
Electrical Behaviour of Condensed
Matter
Objective
Electrical behaviour of solids will be discussed in this chapter. Basis of classifying
solids as insulator, semiconductor, conductor and superconductors will be discussed
in details. After the study of this chapter it is expected that the reader will be able to
understand the behaviour of different crystalline solids when they are subjected to
electric field.
2.1 Introduction
A material may possess several intensive properties that do not depend on the amount
of the material. These quantitative properties are often used as a metric by which the
advantages of one material over the other can be compared for material selection for
a specific purpose. The first and the most important electrical property of a material
is its ‘resistivity’, (or specific resistance) generally denoted by Greek letter ‘ρ’ (rho).
Resistivity is defined as the resistance offered by a unit cube (a block of 1 m × 1 m
× 1 m) of the material between its opposite faces (see Fig. 2.1). The MKS unit of
resistivity is ‘ohm-metre’ written as ‘Ω-m’ in short. The resistance R of a block of a
material of length L and uniform area of cross section A may be written as
L(m)
R(Ω) = ρ(Ω-m) 2
A m
Or

A m2
ρ(Ω-m) = R(Ω) (2.1)
L(m)

https://doi.org/10.1007/978-3-031-32084-2_2
62 2 Electrical Behaviour of Condensed Matter
Fig. 2.1 a Resistance between two opposite faces of a 1 m × 1 m × 1 m cube is equal in magnitude
to the resistivity of the material. b Resistance of a bar of length L and uniform area of cross section
A is given by R = ρ LA
Both resistivity ρ and resistance R are the measure of the opposition that an
electric current faces while passing through the given specimen, however, ρ is an
intensive property of the material (it does not depend on the amount of the material)
while R is an extensive property that depends both on the size and shape of the
material. Resistivity ρ has a fixed value for a given material at a given temperature;
however, for same material it has different values at different temperatures. For most
substances the temperature dependence of resistivity is given as
ρT = ρ0 (1 + K T ) (2.2)
In expression (2.2) ρT is the value of resistivity at absolute temperature T, ρ0 is

the value at absolute zero and K is a constant, called the temperature coefficient of
resistivity. The SI unit of K is Kelvin−1 .
Reciprocal of resistivity is called conductivity and is denoted by Greek alphabet
σ . Conductivity or specific conductance is a measure of the ease with which a current
may pass through the specimen. The SI units of conductivity are reciprocal of the
unit of resistivity and are ‘Siemens per metre’ which in short may be written as
S/m or Sm−1 .

S 1
σ = (2.3)
m ρ(Ω-m)
Magnitudes of the resistivity, conductivity and the coefficient of temperature (K)

for some important materials are given Table 2.1. It may be observed in this table
that the coefficient temperature K has both positive and negative values. For material
where K is positive, the resistivity of the material increases with temperature. The
resistivity of those materials that have negative value of K decreases with the increase
of temperature.
2.1 Introduction 63
Table 2.1 Resistivity, conductivity and temperature coefficient for some materials
Element/material Resistivity (Ω-m) at Conductivity (S/m) at Temperature
20 °C 20 °C coefficient K (K)−1
Gold 2.44 × 10–8 4.10 × 107 3.40 × 10–3
Silver 1.59 × 10–8 6.30 × 107 3.80 × 10–3
Copper 1.68 × 10–8 5.96 × 107 4.00 × 10–3
Iron 9.70 × 10–8 1.00 × 107 5.01 × 10–3
Platinum 1.06 × 10–7 9.43 × 106 3.90 × 10–3
Gallium 1.40 × 10–7 7.10 × 106 4.00 × 10–3
Carbon (amorphous) 5.0 × 10–4 to 8.0 × 1.25 × 103 to 2.0 × − 0.5 × 10–3
10–4 103
Carbon (graphite) 2.5 × 10–6 –5.0 × 2.0 × 105 –3.0 × 105
Parallel to basal plane 10–6 3.30 × 102
Perpendicular to basal 3.0 × 10–3
plane
Gallium arsenide 1.0 × 10–3 –1.0 × 108 1.0 × 10–8 –1.0 × 103
(GaAs)
Germanium 4.60 × 10–1 2.17 − 48.0 × 10–3
Silicon 6.41 × 102 1.56 × 10–3 − 75.0 × 10–3
Diamond 1.00 × 1012 1.00 × 10–13
Teflon 1.0 × 1023 to 1.0 × 1.0 × 10–25 to 1.0 ×
1025 10–23
Elements/materials listed in Table 2.1 may be divided into three main classes;
(a) Have very high or large value of resistivity of the order of ≈ 1025 –1012 (Ω-m).
Such materials are called insulators. Insulators offer very high opposition to
the flow of current. An ideal insulator will have infinite value of resistivity and
will not allow current to pass through it.
(b) Materials that have resistivity value in the range of 102 –10–2 (Ω-m) with nega-
tive value of coefficient K. These materials are called semiconductors and are
extensively used for fabricating electronic devices.
(c) Materials for which the resistivity has small (but non-zero) value ≈ 10–5 –10–10
(Ω-m) are called conductors, and these materials offer small opposition to the
flow of current. It may also be observed that coefficient of temperature K for
conductors has a positive value, meaning thereby that the resistance of a given
piece of the specimen conductor will increase on increasing the temperature.
(d) A fourth category of materials that is not included in the table is called supercon-
ductors. Superconductors show zero resistivity under some special conditions
of temperatures, etc. Having zero resistivity or zero resistance, superconductors
are very important materials as no energy is consumed/wasted in passing current
through them.
In the backdrop of atomic theory of matter, according to which all matter is made
of atoms, the question arises as what is the reason for this difference in electrical
behaviour of different materials? One way of explaining the difference between insu-
lators, semiconductors, conductors and superconductors is in terms of the electron
band theory of condensed matter.
2.2 Electron Energy Band Theory
Most solids are crystalline; they have a regular arrangement of atoms in a pattern
which is repeated in three dimensions. Electric current in solids may flow only by a
net movement of charge carriers that is of electrons, under an electric field. Therefore,
for the flow of current in case of solids it is essential that there are relatively free
electrons that may move when an electric field is applied across it. In contrast, electric
current in liquids may also flow due to the motion of ions under applied electric field,
as it happens in case of electrolytes, etc. In order to understand the physics of current
flow and resistivity in solids, it is required to know: (a) which electrons in the solid
and (b) under what conditions these electrons may become carrier of current. For that
one has to understand the electron configuration in different atoms and how electron
configurations of crystals play a role in current flow.
According to the quantum mechanical model of the atom, electrons of an isolated
atom are distributed in discrete energy levels (or orbitals). Electron energy levels
are discrete but closely spaced. Let us take the example of element Aluminium
each atom of which has a total of 13 electrons. Electron configuration of Al-atom
is 1s2 2s2 2p6 3s2 3p1 , and energy distribution of electrons in different energy levels
for an isolated atom is shown in Fig. 2.2a. Electrons in highest energy shell (3s2 3p1
shell in this case) are called valence electrons. Valence electrons are least tightly
bound with the nucleus of the atom and take part in chemical reactions, etc. If two
atoms of Aluminium come very close to each other and form a diatomic molecule,
the electron energy levels of the di –atomic molecule will be like the one shown in
Fig. 2.2b. There will be two levels close to each other corresponding to the two atoms.
If three atoms come close enough to form a tri-atomic molecule, each level will split
into three closely spaced levels, and so on. In a very small crystal of Aluminium
there is very large number of atoms (≈ 1025 atoms) packed very close to each other,
and therefore, electron energy levels group up in bands of very closely spaced levels
separated by band gaps, as shown in Fig. 2.2c. The energy band that contains valence
electrons is called the valence band and the band next to it in energy the conduction
band. The band energy gap between the valence and the conduction band is called
forbidden energy gap. Energy band gaps are regions of energy where no energy
level of the material exists, and no electron of the atoms of material may have that
energy.
It is obvious that in solids only valence electrons may take part in current flow,
since other inner electrons are tightly bound with the nucleus. On application of an
electric field to the crystal, valence electrons in the crystal are subjected to a force and
2.2 Electron Energy Band Theory 65
Fig. 2.2 a Electron configuration of an isolated atom of Aluminium. b Electron energy levels for
a diatomic molecule of Aluminium. c Electron energy bands in an Aluminium crystal
try to gain energy. Valence electrons will be able to absorb energy only if there are
vacant levels available at higher energies. As such valence electrons will be able to
gain energy and move only if (i) there are vacant levels in the valence band and/or (ii)
when valence band is completely filled then the forbidden energy gap is non-existent
or small enough so that valence electrons may shift to the conduction band, which
is completely empty. Thus current flow in crystalline solids is decided by the nature
of the valence band and the forbidden energy gap. It may, however, be noted that a
partially empty valence band may allow the flow of current but only up to an extent
because of the limitation of available unoccupied levels in the valence band. On the
other hand, if the forbidden energy gap is small or does not exist, then considerable
number of electrons from valence band may take part in current flow as there will be
large number of unoccupied levels available to the electrons in conduction band.
SAQ: No electron level of the parent crystalline solid may exist in forbidden energy
gap, however, if some other atoms are imbedded /mixed in the crystal struc-
ture, electron level corresponding to the other atoms may exist in forbidden
energy gap or not?
So far we made only a qualitative discussion of electron energy bands in a crystal.
Though exact quantum mechanical calculations for many body systems are impos-
sible, however, approximate calculations show that energy gap between two consec-
utive bands depends strongly on the relative separation of atoms. Figure 2.3 shows
the variation of band gap with relative separation of atoms. It is clear from the figure,
Fig. 2.3 Variation of band

gap with relative separation
of atoms
depending on the packing of atom in the crystal the two consecutive bands may be
far apart (as at d1 ), may be separated by a small band gap (d2 ) or may overlap as at
distance d3 .
In the light of the above, three different situations may arise (i) forbidden
energy gap is quite large, and valence shell is completely filled; (ii) valence shell
is completely or partially filled but forbidden energy gap is small and (iii) forbidden
energy gap between the valence and conduction bands is either very small or does not
exist. Corresponding to these conditions materials may be divided into insulators,
semiconductors and conductors.
SAQ: What characteristics of atoms decide the size of forbidden energy gap?
SAQ: Valence and conduction bands in a crystalline solid overlap, the solid will be
a (insulator/conductor/semiconductor)? Choose the correct alternative.
2.3 Insulators
A crystalline solid in which valence shell is completely filled and forbidden energy
gap is quite large becomes an insulator. It is because in such materials valence
electrons do not find vacant levels to move when an electric field is applied to the
specimen. Since only the valence and the conduction bands play role in deciding the
electrical nature of materials, it is customary to show only these bands in pictorial
representation of band structures of solids. The band structure of a typical insulator
is shown in Fig. 2.4a, where the valence band is completely filled with the maximum
number of electrons it can hold and the conduction band is empty, however, the
forbidden energy gap E g between the valence and conduction bands is quite large,
larger than 5 eV. If the band gap is large, electrons do not acquire sufficient energy
from the applied electric field to overcome the band gap and shift to conduction band
2.3 Insulators 67
Fig. 2.4 a Energy band diagram for an insulator b a specimen of length ‘d’ subjected to a voltage
V develops an electric field E = V /d. c Table showing forbidden energy gap for some materials at
specified temperature
where there is large number of vacant levels. Energy acquired by an electron in an

electric field of strength E is eE (eV) See Fig. 2.4b. In case eE < E g , electrons will
not be able to cross the forbidden energy gap, and no current will flow through the
specimen.
It may be noted that in an insulator there are no electrons in conduction band even
at room temperature, and no electrons may jump from valence band into conduction
band by the application of an electric field. Forbidden energy gap is temperature
dependent, and therefore its value is quoted at a given temperature. The strength of
the electric field inside any specimen may be increased by increasing the voltage V.
However, when applied voltage V is increased beyond a given value, called break-
down voltage, the electric field becomes so large that it tears the binding of valence
electrons with atomic core and electrons so released may constitute a current through
the specimen. Hence, insulators may allow some current if the applied voltage is
beyond its breakdown limit. An interesting every day example is the occurrence of
lightning and thunder during raining season. Though air is insulator but under high
electric field produced by oppositely charged clouds, insulation break down of air
takes place producing lightning. The thickness of the insulating material plays a role
in breakdown. Specific dielectric strength is often listed in terms of kilo-volt per
inch (kV/in.) and for some materials listed in Table 2.2.
There is no perfect insulator even the best insulator like ceramic or Teflon may
allow very minute currents even below their breakdown voltage because of impurities
present in their crystals and also because of leakage through surface. Plastic and
rubber insulators that are flexible because of their fiber nature are used to cover
Table 2.2 Dielectric strength

Material Dielectric strength (kV/in.)
of some insulators
Air 50
Porcelain 150
Rubber 600
Paper 1250
Teflon 1500
Glass 2500
Mica 5000
electric wires and cables. Ceramic insulators are generally used for high voltage
applications.
SAQ: What is the physical significance of specific dielectric strength of a material?
2.4 Semiconductors
Those materials for which forbidden energy gap lies in the range of 0.2–3.0 eV and
electron density at Fermi level of around 1020 m−3 are usually classified as common
type semiconductors. These limits are not very rigid some synthetic materials that
have almost 5 eV forbidden energy gap or even larger also behave as semiconductors.
In principle four factors decide the electrical nature of materials, they are (i) the
magnitude of the forbidden energy gap E g , (ii) the magnitudes of the crystal wave
number or crystal momentum vector (k) at the bottom of the conduction band and
the top of the valence band, (iii) number of available electron energy states around
Fermi energy and (iv) the mobility of charge carriers. It may, however, be mentioned
that all these properties are inter-related and are not totally independent.
2.4.1 Intrinsic Semiconductors
Two elements in their very pure form (> 99.999… %) are natural semiconductors;
they are Germanium (Ge) and Silicon (Si). These two naturally occurring semicon-
ductors in their purest form, purity better than 1 part in billion, are called intrinsic
semiconductors. Semiconductors in which some impurity atoms are deliberately
mixed are termed as extrinsic or more frequently doped semiconductors.
It is interesting to observe the difference in electrical properties of elements that
are members of the 14th group of periodic table: 6 C, 14 Si, 32 Ge, 50 Sn and 82 Pb. These
elements have many common properties, for example, all of them have four electrons
in their valence shell, and all have diamond like crystal structure but their electrical
properties are quite different; carbon in diamond form is one of the best insulator,
Silicon is a semiconductor, Germanium is both a semiconductor and half metal, Tin
2.4 Semiconductors 69
(Sn) is a metalloid and lead (Pb) a metal. These differences in electrical behaviour
originate from the difference in the size of their atoms and relative separation between
successive atoms in their crystals. Relative atomic separation in crystalline structure
and the number of valence electrons decides the magnitude of the forbidden energy
gap. In case of diamond the forbidden energy gap is as large as 7 eV, and so it is insu-
lator, while for Silicon and Germanium the forbidden energy gaps are, respectively,
1.02 eV and 0.66 eV, and therefore, the two elements Si and Ge are semiconductors.
Tin has E g of the order of 0.01 eV, and so it is half metal. In the case of lead (Pb) the
valence and conduction bands overlap, and so it is a metal.
SAQ: Which two parameters of atoms of a crystalline solid decide the magnitude
of the forbidden energy gap?
(i) Purification of natural Silicon
The starting material used for the fabrication of semiconductor devices is
monocrystalline Silicon or Germanium. However, most importantly to tech-
nology, Silicon is the principle platform for semiconductor devices. Silicon is
one of the most abundant elements in the crust of earth. The process to transform
raw Silicon into a use-able single crystal substrate for modern semiconductor
processes begins by mining for relatively pure SiO2 . The relatively pure Silicon
dioxide is reduced with carbon in an electric furnace at temperatures ranging
from 1500 to 2000 °C. The reduction process yields metallurgical grade (MG)
Silicon of purity around 97%. However, this Silicon must be further purified to
bring down impurities below the parts-per-billion level. Though several different
methods may be used for further purification of MG-grade Silicon, however,
the two frequently used methods are discussed here.
(a) The trichorosliane method For further purification MG-Silicon is treated
with HCl to form trichlorosilane (TCS) in a fluidized-bed reactor at 300 °C
according to the following chemical reaction
Si + 3HCl → H2 + SiHCl3 (Trichlorosilane TCS) (2.4)
In the process of converting MG-Silicon to TCS, many impurities like

Fe, Al and B get removed. The ultrapure TCS is subsequently vaporised,
diluted with of H2 and flowed into a deposition reactor where it is trans-
formed into elemental Silicon. The contamination level in this Silicon is
typically of the order of 0.001 parts per billion (ppb). It may be further
refined using the zone refining technique.
(b) Zone refining technique
William Pfann, Chemical Engineer who pioneered this technique, repeat-
edly passed a long tube filled with Germanium ingot horizontally through
a series of electric heating coils. This melted portions of Germanium and
allowed them to re-crystallise. The newly crystalline material was found
to be purer than what came before. It was found that impurities become
steadily concentrated in the molten portions, which were swept away. The
technique has undergone several improvements since then.
The present day zone refining technique purifies solids by passing a number
of molten zones through the solid in one direction. Each zone carries a
fraction of the impurities to the end of the solid charge, thereby purifying
the remainder of the solid. The basic lay out of the process is demonstrated
in Fig. 2.5a where the impure crystal of the element (Si or Ge) in the form of
a rod is taken and a heating element that may be moved along the length of
the rod from one end to the other is placed, say, at the extreme left position
indicated by 1. When current is passed through the heating element the
part of the semiconductor rod immediately below the heater melts while
the remaining section of the rod remains in solid state. The heater is then
moved slowly towards right to position 2 and then to 3 and so on. As result
of the motion of heater, successive sections of the semiconductor rod go
to molten state then turn back to solid crystalline states one after the other.
Figure 2.5b shows a typical part of the road, where initially the heater was
above section AB, and the section ABBA of the rod was in molten state.
With time the heater has moved to the location above section BC of the
rod, and now the section of the rod BCCB is slowly turning to molten state.
In the mean time section ABBA, which was in molten state earlier, starts
solidifying, and process of solidification or the process of re-crystalline
starts from the face AA and slowly spreads towards the face BB. Two
important points to note are as follows: (i) the mobility of impurity atoms/
ions is more in molten phase of the semiconductor compared to its solid
crystalline phase, and (ii) the melting point of pure crystalline substance
is always higher than the impure substance. As a result impurities diffuse
from the part of the rod that is undergoing solidification to the part still in
molten phase and pure material crystallises at a higher temperature than
the molten material with impurities. Migration of impurities is shown by
red arrows in Fig. 2.5b. In this way, impurities accumulate at the right side
edge BB of section ABBA. When section BCCB of the rod changes to
Fig. 2.5 a Zone refining technique for improving purity of semiconductors. b Impurities shift from
solid phase to the molten phase
molten state, impurities sitting at face BB migrate to the molten section

BCCB. With the motion of the heating element to the next section of the
road, re-crystallisation in molten section BCCB starts from face BB and
spreads towards face CC. Impurities in section BCCB collect at face CC,
increasing the purity of the re-crystallised solid. Several cycle of motion
of the heating element only in the direction from left to right leaves the
semiconductor ultrapure.
SAQ: Two identical pieces of Silicon semiconductor are given, one is
ultrapure, and the other has high amount of impurities. If both
pieces are heated which one will melt first and why?
(c) Poly crystal to monocrystal
The ultrapure elemental Silicon is, however, poly crystalline and required to
be converted into monocrystalline form. To achieve that, the polycrystalline
Silicon is mechanically broken into 1–3 in. chunks and these chunks then
undergo surface etching and cleaning in a clean room environment. These
chunks are then packed into a quartz crucible for melting in a Czochralski
furnace. A mono crystalline Silicon seed is loaded on seed shaft in the
upper chamber of Czochralski furnace. Slowly the seed crystal is lowered
and dipped up about 2 mm depth in the molten Silicon. The seed is then
withdrawn back up to the surface of the melt, and the melt is allowed to
solidify at the surface. While pulling back the seed, the crucible and the
seed are rotated in opposite directions to grow an almost round crystal.
For proper growth of the monocrystalline crystal, the furnace must be
very stable and free from vibrations. Speed of pulling up the seed and
furnace temperature is two critical parameters. Once the growth process
is complete, the monocrystal is left for cooling inside the furnace for a
considerable period of time up to 7–8 h. The monocrystalline Silicon ingot
so developed has the same orientation as the seed crystal (Fig. 2.6).
In float zone crystal growth technique the end of a long polysilicon rod is
locally melted and brought in contact with a monocrystalline Silicon seed.
The molten zone is slowly moved by moving the heating element (like in
zone refining) leaving behind a ultrapure Silicon monocrystal.
SAQ: Why only ultrapure monocrystal semiconductors are used for
fabricating electronic devices?
SAQ: What is the principle on which the zone refining technique is bases?
(d) Monocrystal to wafers
The monocrystal (also called ingot) is then undergo several precision
mechanical and chemical processes to cut it into wafer of desired size, shape
and other requirement. First of these steps is multiwire slicing (MWS). As
the name suggests, a multiwire slicing machine consists of two cylindrical
rollers with very fine grooves numbering from several hundred to thousand
itched on them. Ultrafine steel wires of diameter around 2 mm are stretched
Fig. 2.6 Sketch of a

Czochralaski furnace
in these roller groves. Both rollers are connected to the same motor and
may rotate at high speed in same direction. Steel wires are painted with
liquid abrasive paint and work as sharp cutting blades. When a monocrystal
is pressed through the wire blades, wafers of Silicon are cut through the
crystal.
A rough sketch of multiwire slicing machine is shown in Fig. 2.7. The
rough drawing is just to understand the principle of working, an actual
multiwire saw is much complicated which has several cylindrical derives
to guide the motion of wire blades and wire spools to maintain a continuous
supply of new wire.
Silicon and Germanium wafers are used for fabricating semiconductor
devices.
(ii) Fermi energy and Fermi level
Electron being a particle having spin 1/2 ℏ is a Fermion and obeys quantum
mechanical statistics called Fermi–Dirac statistics. According to this statistics
at absolute zero temperature (0 K) electrons start filling lowest energy states of
the system (like an atom) obeying exclusion principle, filling in to higher states
after exhausting all lower energy states. The resulting structure of electrons
is termed as ‘Electron Sea’ or ‘Fermi Sea’. The surface of this sea is called
Fermi surface or level. This means at absolute zero temperature no electron can
have energy larger than the Fermi energy Ef of the Fermi level. According to
Fig. 2.7 Sketch of a

multiple wire slicing
machine
Fermi–Dirac statistics the probability ‘p(E)’ of finding an electron with energy

E, at absolute temperature T (Kelvin K) is given as
1
p(E) = (2.5)
1+e ( E−E f )/kT
Here, E f is the Fermi energy for the system and k is Boltzmann constant. Equa-
tion (2.5) gives the theoretical probability of finding an electron with energy E
at temperature T (K). In an actual system there will be an electron with energy
E only if an electron level with energy E actually exists. Often it so happens
that the given system does not have an allowed energy level at some energy,
for example, in case of crystalline solids no energy levels for electrons of their
atoms exist in forbidden energy gap, in that case Eq. (2.5) will still give some
probability of finding an electron for a level in forbidden energy gap.
When one calculates probability of finding an electron with energy equal to
Fermi energy, i.e. if E = E f then Eq. (2.5) gives
1 1 1
p Ef = = = (2.6)
1+e ( E f −E f )/kT 1 + 1 2
Expression (2.6) tells that probability of finding an electron with Fermi energy
E f is 0.5, and that this probability does not depend on temperature. So if
temperature is 0 K or 100 K, the value of p(E f ) will remain 0.5.
Energy band picture for an intrinsic semiconductor at absolute zero and at some
higher temperature T > 0 K is shown in Fig. 2.8.
At T = 0 K, valence band contains all valence electrons, and the conduction
band is empty, which has no electrons. Therefore, the probability of finding an
electron at the bottom of the conduction band is 0, while that of finding an elec-
tron at the top of the valence band is 1. The probability of finding an electron
Fig. 2.8 Band structure of a semiconductor at a Zero Kelvin (0 K) b at temperature T (K) > 0 (K)
with probability p(E f ) = 0.5 will be at a point midway between the bottom
of the conduction band and the top of the valence band, i.e., in the middle of
the forbidden energy gap. Therefore Fermi energy level for an intrinsic semi-
conductor lies in the middle of the forbidden energy gap. Further, since p(E f )
does not depend on temperature, the Fermi level for intrinsic semiconductor will
remain in the middle of the forbidden energy gap at all temperatures. Figure 2.8a
shows the conduction band, valence band and Fermi level for intrinsic semi-
conductor at absolute zero temperature. As may be observed in this figure, all
valence electrons of the intrinsic material at absolute zero temperature are in
valence band, and the conduction band is totally empty.
When temperature of an intrinsic semiconductor is raised above 0 K, electrons
in valence band absorb energy from the surrounding environment and if the
energy gained by a valence electron becomes equal or larger than the forbidden
energy gap E g , the electron jumps to the conduction band leaving an electron
vacancy in the valence band. This vacancy of electron which behaves as a
positive charge is called ‘hole’. Hole is just a fictitious entity but it is very useful
in understanding the physics of semiconductors. The concept of hole originates
from the quantum mechanical treatment of current flow in semiconductors.
It so happens that Schrödinger’s equation when applied to a semiconductor,
under some approximations, separates out into two independent components
one describing the motion of electrons and the other the motion of a positively
charged particle. The absence of electron in the valence band is thus assumed
to be the positive particle, hole, the motion of which is described by the second
component of Schrodinger’s equation. For all practical purposes hole is treated
as a positively charged particle with charge + 1 e and a mass very nearly (slightly
more) equal to the mass of an electron me .
Electrons that have shifted to the conduction band at T > 0 K behave as delo-
calised electrons or free electrons that are not attached to any particular atom
of the semiconductor crystal. Conduction band electrons may be compared to
the delocalised electron cloud in case of metals which is associated with the
crystal lattice but not to any individual atom. When an intrinsic semiconductor
specimen at T > 0 K is subjected to an electric field by applying a voltage
across it, a current consisting of conduction band electrons and valence band
holes may flow through the specimen. Thus at absolute zero a pure or intrinsic
semiconductor behaves as an insulator while at a temperature T > 0 K, the same
specimen behaves as a conductor. Further in a pure semiconductor specimen at
any temperature T > 0 K the number of electrons in conduction band is always
equal to the number of holes in valence band, also the number of electrons in
conduction band (and holes in valence band) increases with the increase in the
temperature T as more valence electrons may jump to the conduction band at
higher temperature.
Fermi level assumes added importance in case of semiconductors, it may be
treated as a reference of energy; energy of electrons in conduction band increases
as one moves upwards from the Fermi level, while the energy of holes increases
as one goes downwards from the Fermi level. That means that an electron at
the top of the conduction band is most energetic while a hole at the bottom of
valence band has largest energy. Further, the probability of finding an electron
say X units of energy above the Fermi level is same as the probability of finding
a hole same X units below the Fermi level.
Figure 2.9 shows the conduction and valence bands for intrinsic Germanium
and intrinsic Silicon crystals at absolute zero and at temperature T > 0 K.
Figure 2.9 is self-explanatory, telling that both intrinsic semiconductors have
empty conduction bands at T = 0 K, and hence behave as insulator. However,
the two points of interest are (i) at T > 0 K, the number of free electrons
in conduction band in Ge is larger than the number of free electrons in Si,
because of the smaller value of its forbidden energy gap. (ii) The energy of free
electrons in conduction band increases vertically upward from the Fermi level,
while energy of holes in the valence band increases vertically downwards from
the Fermi level. Therefore, an electron at the top of the coduction band has
highest energy amongs free electrons, while a hole at the bottom of the valence
band is the one with highest energy amongst holes. Further in a specimen of
an intrinsic material the number of holes is always equal to the number of free
electrons in the conduction band. At any temperature T > 0 K, in an intrinsic
material, new holes and free electrons keep generating on one hand, and on the
other hand they also get annihiliated when a free electron falls back in valence
band and recombines with a hole. The process of electron + hole generation
and annihilation goes on simultaneously in such a way that the average number
of holes and free celectrons remain constant over a period of time.
SAQ: What is Fermi level and what is its physical significance?
Fig. 2.9 a Intrensic Ge crystal at 0 K. b Intrensic Si crystal at 0 K. c Intrensic Ge crystal at T >

0 K. d Intrensic Si crystal at T > 0 K
SAQ: Do holes may move in intra atomic space like free electrons? Justify
your answer.
2.4.2 Covalent Band Picture of Intrinsic Semiconductor
In the previous section we studied the electron band theory and its application in
distingushing different types of crystalline materials according to their electrical
properties. An other equivalent way of describing electrical properties of crystalline
solids is by using the covalent bond picture of these materials. Atoms of Silicon and
Germanium in their intrinsic crystals are held together by covalent bonds. Both these
atoms (Ge and Si) have four valence electrons in their valence shells. Electronic
configuration and arrangnment of electrons in a 14 Si atom are shown in Fig. 2.10.
Since only the valence electron of an atom take part in bonding and inner electrons
plus nucleus does not play any role, it is convinient to represent an atom by a core
(that has nucleus and inner electrons) with positive charge equal to the number of
valence electrons, (+ 4e) in case of Silicon, and four valence electrons. As a matter of
fact any atom with four valence electrons (like Si, Ge, Sn, Pb) may be represented by
a core of + 4e charge and four electrons. An atom for example of Aluminium which
has three valence electrons may be represented by a positively charged core having
Fig. 2.10 Representing an

atom of 14 Si with four
valence electrons by a core
of + 4e charge surrounded
by four valence electrons
charge + 3e and three valence electrons and a penta valent atom (like phosphorus)
by a core of + 5e charge and five electrons.
In covalent bonding atoms share their electrons; in case of Si and Ge four neigh-
bouring atoms share their one electron each forming the crystal lattice. Figure 2.11
shows the covalent bonding in Silicon or Germanium intrinsic crystal at absolute 0 K
temperature.
As may be seen in the figure, at T = 0 K all covalent bonds are intact and all valence
electrons of each atom are shared by neighbouring four atoms. Since electrons are
held in covalent bonds because of the force of attraction of nearby positively charged
Fig. 2.11 Pictorial representation of covalent bonding in intrinsic Silicon or Germanium crystal at
absolute zero temperature
cores, they cannot move even when an electric field of moderate strength is applied
to the crystal, hence the crystal behaves as an insulator at T = 0 K.
Covalent bonds are characterised by bond energy; the energy by which electrons
are held within the bond. The covalent bond energy of Silicon is 1.08 eV and that of
Germanium 0.66 eV. It means that if an electron in the covalent bond of Silicon crystal
somehow gets energy either equal to 1.08 eV or large; it may break the covalent bond
and will become a delocalized or free electron which may hop from one atom to the
other in the intra atomic space of the crystal. It may be observed that the covalent
bond energy is simply equal to the forbidden energy gap of band theory. Electrons
may get energy in several ways, like by heating or by putting them in an electric field
etc.
If an intrinsic Silicon or Germanium crystal is heated to say some temperature
T > 0 K, its electrons in covalent bonds will acquire temperature T > 0 K and will
get thermal energy ≈ kT, where k is Boltzmann constant (k = 8.6 × 10–5 eV per
Kelvin). The most likely form of this thermal energy is kinetic energy associated
with vibrations, electrons which were stationary in covalent bonds at T = 0 K, starts
vibrating when the temperature of the specimen is increased. When thermal energy
of the electron increases beyond the bond energy (1.08 eV for Si and 0.66 eV for Ge),
it may break the bond and come out of the bond becoming a delocalised electron,
leaving a hole in the covalent bond. The structure of an intrinsic semiconductor
crystal (Si or Ge) at T > 0 K is shown in Fig. 2.12
The process of bond breaking resulting in creation of electron–hole pairs and
the opposite process of recombination of electron–hole pairs to remake some of
the broken bonds simultaneously keep going in the crystal at temperature T > 0 K.
Ultimately, equilibrium is reached when the rate of new electron–hole pair creation
becomes equal to the rate of recombination of electron–hole pairs.
Covalent bond breaking →→ free electron + hole (2.7)
Free electron + hole →→ reformation of covalent bond (2.8)
At equilibrium, rate of reaction given by Eq. (2.7) becomes equal to the rate of
process represented by Eq. (2.8). As a result at any temperature T > 0 K; the average
number of free electrons and holes per unit volume of the semiconductor becomes
constant (Fig. 2.12).
Further, at equilibrium the average number per unit volume (called the number
density or carrier concentration) of free electrons and holes in the crystal is equal.
The average number density of electrons and holes in an intrinsic semiconductor
is equal and constant at a fixed temperature; however the average number density
increases with the rise of temperature. Further, it is a common practice to call free
or delocalised electrons simply as electrons and instead of saying average number
density, simply to say number density. Though obvious, but one must remember that
free electrons (these are the electrons which in electron band theory shifts to the
conduction band on acquiring of energy) are free to move within crystal while holes
Fig. 2.12 Breaking of covalent bonds in intrinsic semiconductor at T > 0 K
are always bound within the covalent bond and can move from one covalent bond to
the next.
If the (average) number densities of (delocalised) electrons and holes in an intrinsic
semiconductor at temperature T are, respectively, denoted by n ie and n ih , then
n ie = n ih (2.9)
The value of free electron number density n ie for Silicon at 300 K (nearly room
temperature) is of the order of ≈ 1.08 × 1010 cm−3 .
When a specimen of intrinsic semiconductor at T > 0 K is subjected to an electric
field by applying a voltage across its two opposite faces, the free electrons in the
specimen moves within the enter atomic space of the crystal in a direction opposite
to the electric field, as shown in Fig. 2.13. Holes that are bound in covalent bonds
also shift from one covalent bond to the next in the direction of the imposed electric
field (see Fig. 2.13). Resultant current I in the circuit is the sum of the electron and
hole currents. Therefore, both the electrons and holes participate in current flow in a
semiconductor.
SAQ: Consider Fig. 2.13 and draw three successive steps showing the motion of a
hole under applied voltage V.
2.4.3 Doped or Extrinsic Semiconductors
Addition of a very small amount (of the order of 1 part in 107 parts) of impurities in
an intrinsic semiconductor crystal may drastically change its electrical conductivity,
Fig. 2.13 Motion of free electrons and holes under electric field
Fig. 2.14 Block diagram of hybrid medium current ion implanter
optical and structural properties. The crystal with controlled impurity added to it is
called a doped or extrinsic semiconductor crystal. The process using which small
amount of impurities is added in a controlled way is called doping.
The reason why deliberately added impurities are mixed only in small amount is
(i) to avoid any breakdown in crystal structure and (ii) addition of very small amount
of impurity is sufficient to change the conductivity of the intrinsic semiconductor to
the desired value. Addition of excessive impurity may turn a semiconductor into a
conductor. Further, only two types of impurities are added; impurity atoms with either
5 valence electrons (pentavalent) or with 3 valence electrons (trivalent). Pentavalent

atoms that are often used as dopant are arsenic (As), phosphors (P) and Antimony
(Sb). Similarly, trivalent impurities that are often used for doping are boron (B),
Aluminium (Al), gallium (Ga) and indium (In). Pentavalent impurities are called
donor impurities while atoms of trivalent impurities are termed as acceptor impurity.
2.4.4 Doping Technology
Doping refers to the process of introducing impurity atoms in a semiconductor in a

controllable manner in order to define the electrical properties of the semiconductor.
Doping of an intrinsic semiconductor by donor (pentavalent impurity atoms) and
acceptor (trivalent atoms) atoms may change the free electron and hole concentra-
tions (number density) in the doped semiconductor from 1013 cm−3 to 1021 cm−3 .
Controlled doping may also change the spatial distribution of carriers (free elec-
trons and holes) in the specimen semiconductor quite accurately. Spatial variation
in carrier concentration is required in fabricating devices like pn junction diode and
transistors.
(i) Ion implantation technology
Ion implantation is the front line technology of the day for semiconductor
doping. In this technology an accelerated beam of the dopant ion (like that
of boron ion or phosphorous ion etc.) of controlled energy and flux is made
to hit the semiconductor wafer/monocrystal. Depending on the energy of the
beam, the incident dopant ions penetrate up to a certain depth in the target
wafer/monocrystal. The number density of dopant atoms in the target wafer
may be controlled by controlling the flux and time of irradiation by dopant.
The incident dopant ion beam may be masked externally by placing sufficiently
thick shields in the path of the ion beam; internally within the semiconductor
channelling may take place by crystal structure that is not transparent to the ion
beam. Internal channelling is called self-aligned implant. Main reason for the
popularity of ion implantation technique is the accuracy with which dopants
may be distributed in the target semiconductor. Ion implant technology got
a big boost with the development of high current ion accelerators. Most of
commercial ion implanters have a linear accelerator which may accelerate ions
up to several MeV of energies. Ion source, which produces the desired ions
and accelerating column, where ions are accelerated to higher energies, is the
two main components of an accelerator. Earlier accelerators used cold cathode
ion source that could deliver ion currents of only few hundred micro appears
(10–6 A). However, in the year 1970 a new type of ion source called hot cathode
ion source was developed that may deliver ion currents of few hundred milli
appears (10–3 A) at around 80 keV ion energy. Today, different ion implanter
machines are available that cover the entire range of both energies and beam
current requirement for semiconductor fabrication industry. These machines

may be grouped as medium current, high current and special implanters.
Most of the medium current implanters that deliver ion beam currents of few
mA use the concept of hybrid scanning. These systems have a hot cathode ion
source where the ion beam is generated. The ion beam then passes through
a system called mass analyser. Mass analyser separates out ions of different
masses and focuses them at different points in its focal plane. Any undesired
ion that may be present in the ion source may be rejected at this stage. The
beam of selected desired ions is then accelerated/decelerated in an accelerating
column. The beam coming out of the accelerating column is made to go through
an energy filter that ensures that ions of only the desired energy may travel
further. The energy filtered ion beam then passes through a scanning magnet
that scans and collimates the ion beam. In one-dimensional scan, the scan
magnet scans the collimated monoenergetic and highly pure ion beam in one
direction say in horizontal direction. That means that the ion beam is slowly
swept in horizontal direction at a certain sweep rate. The target holder which
holds the semiconductor wafer is scanned in vertical direction at the same rate.
As a result of this hybrid scan, the sample wafer is uniformly irradiated by the
ion beam. Block diagram of a hybrid scan medium current ion implanter is
shown in Fig. 2.14.
In normal case ion beams of energies ranging from about 100 eV to 3 MeV
and irradiation dose (number of incident ion per unit area) ranging from 1011
to 1016 ions cm−2 are used.
After ion implantation the crystal structure is damaged by incident ions. Further,
the dopant lions are electrically inactive as majority of them do not occupy
positions in crystal lattice. The ion implanted wafer/crystal is then subjected
to thermal annealing process that restores the damaged crystal structure and
imbed implanted ions in crystal lattice turning then electrically active.
Figure 2.15 shows the profile of implanted dopant ions inside the semicon-
ductor wafer. Sides of the wafer are covered by mask material that does not
allow incident ions to enter the semiconductor and define the volume of the
wafer exposed to implantation. Since charged ions have a certain range in semi-
conductor material, the maximum deposition of dopant ions takes place at a
depth equal to the range of the incident ion. Range depends on the energy of
the ion and the material of semiconductor. Therefore, by selecting the proper
energy of incident ion beam the depth of maximum deposition of dopant may
be adjusted.
SAQ: Explain why the impurity ion concentration shows a peak at certain
depth in the semiconductor when impurity is introduced by ion-
implanting method?
(ii) Diffusion technology
Fig. 2.15 a Process of ion implantation. b Dopant profile in the ion implanted semiconductor wafer.
Highest dopant concentration at range equivalent depth
Diffusion is another important method of introducing dopant atoms in semicon-

ductor material in a controlled way. In diffusion process dopant atoms are intro-
duced in the wafer from gas phase. The diffusion process may be considered
as a series of atomic movement of the dopant atom in crystal lattice.
Figure 2.16 shows the two mechanisms of diffusion mechanism. When temper-
ature of a semiconductor wafer/monocrystal is raised, the host atoms (atoms
of semiconductor material) start vibrating. When kinetic energy gained by a
particular host atom accede its binding energy with crystal lattice, the atom
leaves the lattice and moves to the interstitial space. Thus a vacancy is created
in crystal lattice that may be occupied by the dopant atom (see Fig. 2.16a)
which are in gaseous phase. The other method of doping via diffusion is that
initially the dopant atom diffuses in interstitial space, as shown in Fig. 2.16b
and later on occupies a site in crystal lattice when on further heating a vacancy
is created.
Diffusion is the process in which atoms/molecules moves from region of higher
concentration to the region of lower concentration. Therefore, concentration
gradient is the driving force for diffusion. When a certain area of the surface
Fig. 2.16 a Vacancy mechanism of diffusion. b Interstitial mechanism of diffusion

Fig. 2.17 a Process of diffusion. b Profile of diffused atoms
of a semiconductor wafer is exposed to dopant atoms in gaseous phase, dopant

atoms diffuse from the region of higher concentration (container of gaseous
dopant atoms) into wafer, where initially there were no dopant atoms (see
Fig. 2.17a). The flux F (number of atoms per unit area per unit time).
Is proportional to the gradient of concentration of dopant atoms, i.e. for a
one-dimensional case
dC dC
F ∝− or F = −D (2.10)
dx dx
Here, dC
dx
is the concentration gradient of dopant atoms in direction x and
the negative sign signifies that motion of diffusing atoms is in the direction
of higher to lower concentration. The constant of proportionality D is called
diffusion coefficient. With diffusion of dopant from gaseous phase container
into semiconductor wafer, the concentration difference on the two sides will
decrease and assuming D to be constant, and flux of diffusing atoms will also
decrease. Ultimately, diffusion of dopant atoms will stop when concentration
of atoms in container and in wafer will become equal. The depth of diffusion
in wafer essentially depends on the temperature of dopant atoms, higher the
temperature larger the diffusion depth. Depth profile of diffused atoms in the
wafer at a given temperature is shown in Fig. 2.17b.
According to the law of conservation of matter, the change of dopant concen-
tration with time must be equivalent to the local decrease of the diffusion flux,
in the absence of a source or sink, therefore,

∂C ∂F ∂ ∂C ∂C ∂ 2C
=− = D or =D 2 (2.11)
∂t ∂x ∂x ∂x ∂t ∂x
Equation (2.11) is often called Fick’s Second law of diffusion.

Boron is the most common trivalent impurity dopant in Silicon; whereas arsenic
and phosphorous are used extensively as pentavalent impurities. These three
elements are highly soluble in Silicon with solubilises exceeding 5 × 1020
atoms per cc in the diffusion temperature range between 800 and 1200 °C.
These dopants can be introduced via several means, including solid sources
(BN for B, As2 O3 for As and P2 O5 for P), liquid sources (BBr for B, AsCl for
As and POCl3 for P) and gaseous sources as B2 H6 , AsH3 and PH3, respectively,
for boron, arsenic and phosphorous.
SAQ: Two exactly identical wafers of Silicon are doped using diffusion tech-
nology. The temperatures of diffusing gas are different for the two
wafers. What will be the difference in the profile of doped impurity in
two cases?
(iii) Doping at monocrystal growth stage
Impurities in controlled amount may also be introduced in a semiconductor
monocrystal while it is being grown from poly crystalline melt. If calcu-
lated amount of either tri or pentavalent dopants is mixed with the molten
polycrystalline a part of these impurities enter the monocrystal.
2.4.5 n and p Type Semiconductors
(i) n-type semiconductor

Doping an intrinsic semiconductor wafer with pentavalent impurity, like phos-
phorus, arsenic or Antimony, results in n-type semiconductor. As a result of
controlled doping the pentavalent impurity atoms occupy regular places in
crystal lattice. Since the number of doping atoms is much smaller than the
number of intrinsic atoms (Ge or Si), the impurity pentavalent atoms are
surrounded by atoms of intrinsic material on four sides. The impurity atom
is surrounded by four covalent bonds of neighbouring atoms and four out of
five valence electrons of impurity atom are held in these four covalent bonds.
However, the fifth electron of impurity atom remains loosely attached with the
parent atom. As a result at absolute zero temperature the n-type material behaves
like an insulator, just like intrinsic material.
At temperature T > 0 K, the loosely bound fifth electron of pentavalent impurity
atom becomes free, leaves the parent atom and like other conduction band
electrons may take part in current flow. With the loss of fifth electron the impurity
atom gets ionised and becomes a positive ion. The covalent bond picture of n-
type semiconductor is shown in Fig. 2.18. At any temperature T > 0 K, each
impurity atom gives or donates one additional free electron to the conduction
band, therefore, the pentavalent impurity is called donor impurity. Since the
concentration of donor impurity is of the order of 107 impurity atoms per cubic
centimetre, the number of additional free electrons provided by donor impurity
Fig. 2.18 Covalent bond structure of an n-type semiconductor
is of the order of 107 free electrons per cm3 . Apart from free electrons given by
impurity atoms, there are also some electron and holes that are created by the
breaking of covalent bonds. But the number of electrons given by donor impurity
per unit volume is much larger than the number of free electrons produced by
covalent bond breaking. Thus, at T > 0 K, an n-type semiconductor has (a) large
number of free electrons due to the ionisation of donor atoms + (b) large number
of positive ions of donor atoms that are fixed in crystal lattice and cannot move
+ (c) few electrons due to bond breaking and + (d) few holes due to covalent
bond breaking. Both free electrons and holes take part in current flow if an
n-type semiconductor is subjected to electric field by applying voltage across
it. Since the number of free electrons in n-type semiconductor is much larger
than the number of holes, free electrons are called majority carriers and holes
the minority carriers.
Electron band diagram of an n-type semiconductor is shown in Fig. 2.19ii where
the band picture for an intrinsic semiconductor is also given for comparison.
Three major differences between the two diagrams may be observed in Fig. 2.19;
(a) the conduction band of n-type semiconductor has large number of free elec-
trons, most of which are due to the ionisation of donor impurity atoms. (b) There
is a donor level, an energy state for impurity atoms, just below the bottom of
conduction band in forbidden energy gap. Energy difference between the bottom
of conduction band and donor level is equal to the binding energy of the fifth
electron of donor atom which could not be accommodated in covalent bonds. It
sometimes appear confusing that how there could be any level or energy state
in forbidden energy gap? At this point it is important to realise that forbidden
energy gap region is restricted only for energy states or levels of parent intrinsic
atoms, atoms of Si or Ge cannot have energy levels in forbidden energy gap;
however, other impurity atoms may have their energy states in the region of
forbidden energy. Further, Fig. 2.19 refers to a temperature T > 0 K, the temper-
ature at which most of the impurity donor atoms have ionised and thus one
electron from each donor atom has jumped to the conduction band leaving a
positive donor ion at donor level. Positive donor ions in Fig. 2.19 are repre-
sented by the + sign. (c) An other very crucial observation from the figure is
that Fermi level of n-type material has shifted up words, towards the conduction
band, from the middle of forbidden energy gap. This has happened because of
the large number of electrons in conduction band as compared to the holes in
valence band. In intrinsic material at any temperature T > 0 K, the number of
electrons in conduction band is exactly equal to number of holes in valence
band as both are generated by the process thermal braking of covalent bonds.
An n-type semiconductor contains large number of free electrons as majority
carriers, few holes from covalent breaking, as minority carriers (see Fig. 2.20)
and large number of immobile positive ions of donor atoms fixed in crystal
lattice.
Fig. 2.19 Electron band diagram of i an intrinsic semiconductor at T > 0 K ii an n-type

semiconductor at T > 0 K
Fig. 2.20 Pictorial

representation of n-type
semiconductor
(ii) p-type semiconductor
Doping of an intrinsic semiconductor by some trivalent impurity, like boron,

Aluminium, or indium results in a p-type semiconductor. A trivalent impurity
atom when trapped by tetravalent intrinsic atoms on all four sides develops
four covalent bonds with four neighbouring atoms for the continuity of lattice
structure. However, one of these covalent bonds has an electron vacancy as
trivalent atom has only three electrons to shear in four covalent bonds. This
electron vacancy is nothing but hole. Therefore each trivalent impurity atom
has one hole in one of the covalent bonds associated with it. At absolute zero
temperature all electrons in covalent bonds are fixed at their locations as they
have no energy. When temperature of the p-type material is raised above 0 K,
electrons may move within covalent bonds and by chance an electron from
some other nearby covalent bond may drop in the vacant place, the hole asso-
ciated with trivalent impurity atom. This impurity atom will then have all the
four covalent bonds around it filled with electrons, and therefore the trivalent
impurity atom will become a negative ion, while hole will be shifted in some
other covalent bond. Thus at T > 0 K holes associated with trivalent impurity
atoms become mobile and each impurity atom becomes a negative immobile ion
fixed in crystal lattice. Since the trivalent impurity atoms accept an electron and
add one additional hole per impurity atom, they are called accepter impurity.
Covalent bond picture of a p-type semiconductor is given in Fig. 2.21.
A p-type semiconductor at T > 0 K has large number of mobile holes which
play the role of majority carriers while few free electrons produced by bond
breaking work as minority carrier.
Energy band diagram of an intrinsic and a p-type semiconductor at T > 0 K are
shown in Fig. 1.22. It may be observed in this figure that for p-type material (i)
conduction band has few free electrons essentially from covalent bond breaking
Fig. 2.21 Covalent bond picture of a p-type semiconductor

Fig. 2.22 Energy band picture at T > 0 K for a intrinsic semiconductor. b p-type semiconductor
that are minority carriers; (ii) the Fermi level is below the middle of forbidden
energy gap towards the valence band; (c) there is an acceptor level just above
the top of valence band holding negative ions of impurity acceptor atoms; (d)
there are large number of holes in valence band most of which are created by
the ionisation of trivalent impurity acceptor atoms (Fig. 2.22).
As in case of n-type semiconductor, the semiconductor is over all neutral,
similarly, p-type material is also over all neutral.
As shown in Fig. 2.23, a p-type semiconductor has holes as majority carriers,
electrons produced by bond breaking as minority carriers and negative ions of
acceptor impurity atoms fixed in crystal lattice as immobile charges.
The majority charge carriers, electrons in case of n-type semiconductor and
holes in p-type semiconductor are released by the impurity atoms only when
Fig. 2.23 Pictorial

representation of p-type
semiconductor
these atoms get ionised. Those impurity atoms that get ionised at room tempera-
ture and release majority carriers are called shallow impurities. The ionisation
energy for shallow impurities is of the order of ≈ kT where k is Boltzmann
constant and T ≈ 300 K.
2.4.6 Compensated Semiconductor
When an intrinsic semiconductor is simultaneously doped by both the shallow donor

and shallow acceptor impurities it is termed as compensated semiconductor. If N D
and N A , respectively, denote the concentration of donor and acceptor impurities, when
all impurity atoms have ionised, and N D = N A , then the compensated semiconductor
behaves as an intrinsic semiconductor. However if N D > N A , compensated material
behaves as n-type and if N A > N D then as p-type semiconductor.
SAQ: In order to fabricate n-type and p-type semiconductors the intrinsic material
is doped using compensated method; can you give a reason why simultaneous
doping with donor and acceptor impurities is better than doping by a single
type of impurity?
2.4.7 Degenerate and Non-degenerate Semiconductors
When an intrinsic material is doped with small amount of impurity, the impurity
atoms in crystal lattice are far apart and do not interact with each other. Such doped
semiconductors are termed as non-degenerate semiconductors. The impurity level
(donor level in case of n-type and acceptor level in p-type) in non-degenerate semi-
conductor is discrete and sharp. However, if impurity concentration is relatively high
and impurity atoms in crystal lattice are near to each other, they interact and the impu-
rity level does not remain a single discrete level but it becomes a band of many energy
levels and sometimes this impurity band may overlap with the nearby (conduction
band in case of n-type and valence band in case of p-type semiconductor) band of
the semiconductor. Such semiconductors are called degenerate semiconductors.
Most of the doped semiconductors that are frequently used are non-degenerate
type.
SAQ: Draw a rough sketch of energy band picture for a degenerate p-type material.
2.4.8 Direct and Indirect Semiconductor
The free electrons in the conduction band of a semiconductor are not really free; their
motion is constrained by the periodic potential of ions in crystal lattice. The effect of
periodic potential is included by assigning an effective mass to electron, denoted by
m ∗n . Similarly, the motion of holes in covalent bonds is also restricted, and therefore,
holes are also assigned an effective mass m ∗h . The energy momentum relation for a
conduction band electron may be written as
pc2
E=
2m ∗n
Here pc is crystal momentum along a given crystal direction defined by the Miller
index. In some semiconductor materials, like Silicon (Si) and gallium arsenide
(GaAs), the maximum of valence band energy and the minimum of the conduction
band energy lie along the direction defined by pc = 0. As a result in such semiconduc-
tors a transition across forbidden energy gap requires just the absorption or emission
of energy. Those semiconductors for which the maximum of valence band energy
and the minimum of conduction band energy lie in the same crystal direction or have
same value of pc are called direct semiconductors. The time response of direct semi-
conductors to photon absorption/emission (Si and GaAs) is fast, and therefore they
are frequently used in optoelectronic devices. On the other hand, those semiconduc-
tors in which the maximum of valence band energy and the minimum of conduction
band energy have different values of crystal momentum pc are called indirect semi-
conductors. Transition across forbidden energy gap in indirect semiconductors is
slow because of the different values crystal momentum.
SAQ: What may happen when a semiconductor absorbs a photon?
SAQ: The photon absorption and emission processes in direct semiconductors are
fast; why?
2.4.9 Compound Semiconductors
Several compounds of trivalent and pentavalent elements, like gallium nitrate

(GaN), gallium arsenide (GaAs), indium phosphate (InP), etc., behave as semicon-
ductor. Most of these compound semiconductors are direct semiconductors and have
forbidden energy gaps in the range of 1.4–1.5 eV. These compound semiconductors
are, therefore, very suitable for fabricating optoelectronic devices. A large band gap
reduces background noise and makes devices stable. Covalent band structure of a
compound semiconductor (InP) is shown in Fig. 2.24. The total number of valence
electrons in the two atoms (3 + 5) is eight which distribute themselves in four covalent
bands as shown in the figure.
SAQ: Give a reason why GaAs is often used for fabricating optoelectronic devices.
Fig. 2.24 Covalent band structure of a compound semiconductor
2.4.10 Current Flow in Semiconductor
Both intrinsic and semiconductors doped with shallow impurities have free electrons
and holes at room temperature. Paul Drude, a German physicist, argued that free
electrons in semiconductors and also in conductors (metals), under some assumptions
may be treated as molecules of an ideal gas. He argued that kinetic theory of gases
which in principle is applicable to an ideal gas may also be applied to free electrons.
The law of equipartition of energy, which follows from kinetic theory, says that
1/2 kT of energy is associated with each degree of freedom, a molecule of an ideal
gas that has three degrees of freedom, at temperature T > 0 K possesses 3/2 kT
of kinetic energy. As such a free electron in semiconductor at room temperature
≈ 300 K will have roughly 6.21 × 10–21 J as kinetic energy. If Vthe is the speed of
electron at room temperature and 9.1 × 10–31 kg the mass of electron, then
1
× 9.1 × 10−31 × Vthe
2
= 6.21 × 10−21 (2.12)
2
Solution of the above equation gives V the , the thermal velocity of a free electron
at room temperature, to be of the order of 1.17 × 105 m/s. This means that at
room temperature free electrons in a semiconductor are moving with high speed of
the order of 105 m/s. If no external potential is applied to a semiconductor and if
there is no charge gradient within the semiconductor, then free electrons in it will
be moving in random directions with velocities of 105 m/s. Fast moving electrons
frequently collide with vibrating crystal lattice (crystal lattice also vibrates because
of temperature), and at each collision its direction of motion gets changed. Hence
in absence of any external voltage and any static charge gradient, the net effect of
lattice-electron collisions is that on average equal number of electrons move in all

directions at a given instant, and hence there is no net current in any direction.
Current may, however, be made to flow in a semiconductor piece by many
processes. Some of the important processes that may cause current to flow in a
semiconductor are (a) application of external voltage that generates an electric field
in semiconductor, (b) accumulation of charge at some location that produces static
charge gradient within semiconductor, (c) injecting additional excess charges, (d)
high field operation. We shall, however, discuss in brief only first two mechanisms
of current flow in a semiconductor.
(i) Drift current Fig. 2.25 shows apiece of semiconductor of length D to which
an external voltage V is applied to establish an electric field E = VD within the
semiconductor in X direction.
Each free electron in semiconductor experiences a force F = −(eE) in direction
opposite to the direction of electric field. Force F produces acceleration a =
F/m ∗n in the motion of electron. Here m ∗n is the effective mass of the electron.
If τ is the time between two successive collisions of the electron with vibrating
crystal lattice and if it is assumed that electron comes to a momentary rest after
each collision, then the velocity gained by the electron in direction opposite to
the electric field v nD , called drift velocity, is given by
e.τ e.τ
v nD = a.τ = − ∗
ε = −μe ε, where μe = ∗ (2.13)
mn mn
Drift velocity, velocity acquired by an electron in an external electric field, is

proportional to the strength of the electric field, and the proportionality constant
μe is called electron mobility. Electron mobility depends on collision time τ ,
which in turn depends on temperature. Similarly one may define hole mobility
μh = me.τ∗ and hole drift velocity v hD = μh ε.
p
Drift velocity v nD gets superimposed on each electron along with the thermal
velocity Vthe under the influence of the electric field. It is easy to show that
Fig. 2.25 Force experienced

by an electron in an electric
field
electron current density due to drift velocity j De is given by
j De = −e.n e .v nD = e.n e μe .ε (2.14)
Similarly, the hole current density due to drift is given as
j Dh = e.n h .v hD = e.n h μh .ε (2.15)
And the total current density due to drift j D is given as
j D = j De + j Dh = eε[n e μe + n h μh ] (2.16)
ne , nh , μe , and μh in above expressions are, respectively, the concentration of

electrons, concentration of holes, mobility of electrons and mobility of holes.
But,
J D = σ ε, hence, σ = e[n e μe + n h μh ] (2.17)
(ii) Diffusion current Diffusion is a universal phenomenon in which particles or

in general objects, irrespective of their charge, move out from a region of
higher concentration to the region of lower concentration so as to equalise
the concentrations on the two sides. If in a semiconductor there is concentration
gradient of charge carriers, electrons and holes, the carriers will move from
higher concentration side to the lower concentration side establishing diffu-
sion currents. Charge carrier concentration gradient in a semiconductor may
be produced by non-uniform doping or by injecting charge carriers of a given
type by implantation.
Let us consider the case where in a semiconductor piece there is electron concen-
tration gradient dn
dx
e
in x direction. In absence of any electric field the free elec-
trons will be moving in random directions with thermal velocity Vthe . Electrons
will undergo collisions with vibrating crystal lattice every now and then and if
τ is the average time between two successive collisions, then one may assign a
mean free path λ = V the . τ . Under these conditions it can be shown that electron
diffusion current density jeDiff and hole diffusion current density jhDiff are given
by
dn e dn h
jeDiff = e.Vthe .λ. and jhDiff = −e.Vthe .λ. (2.18)
dx dx
Putting Vthe .λ ≡ Dn for electron and Vthe .λ ≡ D p for holes, Eq. (2.18) reduces
to
dn e
jeDiff = e.Dn . (2.19)
dx
And
dn h
jhDiff = −e.D p . (2.20a)
dx
Dn and Dp in above expressions are called diffusion coefficient, respectively,

for electron and hole. Though it may appear surprising at first sight but diffusion
coefficient is related to the corresponding mobility. It can be shown that
kT kT
Dn = μe and D p = μ p (2.20b)
e e
2.4.11 Temperature Dependence of Semiconductor Resistivity
As already discussed, resistivity of a material is a measure of the opposition offered

to the flow of current by a unit cube of the material. If a voltage V is applied across
the two opposite faces of a cube of a given material and a current I passes through
the circuit, then resistivity ρ = V /I. For a fixed value of V, resistivity is inversely
proportional to current I. At a given temperature, the value of current I depends
on three factors: (i) the number density or concentration of charge carriers (both
free electrons and holes are charge carriers in case of semiconductors while only
free electrons are charge carriers in conductors), (ii) number of collisions per unit
time between charge carriers and crystal lattice and (iii) trapping of charge carriers
per unit time at trap locations. Higher concentration of charge carrier results in
increased current and in turn decreases the resistivity. On the other hand motion of
charge carriers under the influence of the applied electric field gets randomised due
to charge carrier-lattice collisions which reduces current and increases resistivity. It
may be said that collisions between charge carriers and crystal lattice are the most
frequent and important as collisions between charge carriers themselves are very
less probable on account of their extremely small size and relatively high speed of
their motion. Even in a very pure crystal there are always some sites where charge
carriers are trapped. Trapping of charge carriers reduces current in turn increasing the
resistivity. However, the probability of charge carrier trapping does not very much
depend on temperature.
At room temperature (≈ 300 K), in a semiconductor doped with shallow impurity
of the order of 1013 atoms per m3 (107 atom/cm3 ) there are around 1013 majority
charge carriers per metre cube and some 104 –105 minority carriers per cubic metre
due to covalent bond breaking. There is also a certain rate of lattice-carrier collision,
say N c collisions per second. On raising the temperature above room temperature,
more covalent bonds got broken and more charge carriers are produced. Generation
of additional charge carriers increases almost exponentially with temperature. On the

other hand, collision rate N c also increases but linearly with temperature. Thus, rate
of generation of additional charge carriers wins over the increase in collision rate N c .
The net result of temperature increase is that current I through the semiconductor
increases with temperature, reducing the resistivity of the semiconductor.
In conductors (metals) at room temperature all atoms of the conductor are ionised
and contribute at least one free electron, therefore, the free electron concentration in
conductors is of the order of 1028 free electrons per m3 . As such the resistivity of
conductors is very small of the order of 10–8 (Ω-m) (see Table 2.1). When temper-
ature is raised, in case of conductors there is no further increase in the number of
charge carrier but the rate of electron–lattice collision increases. As a result current
I decreases with temperature, increasing the resistivity of conductors.
SAQ: The mobility (μ) and the coefficient of diffusion (D) are parameters of two
totally independent causes of current flow in a semiconductor but these
constants are related with each other. What is inter-connecting link of the
two?
2.4.12 Theoretical Calculation of Carrier Concentration

in a Semiconductor
There are two types of charge carriers, such as free electrons in conduction band
and holes in valence band in a semiconductor. The concentration or number of
charge carriers per unit volume of the semiconductor which is also called the carrier
number density depends on the concentration of dopant impurity and temperature.
Dopant concentration essentially decides the concentration of majority carriers while
temperature, that determines the rate of covalent bond breaking, controls the minority
carrier concentration.
Both conduction band and the valence band have large number of discrete energy
states where carriers may reside. These energy states, though discrete, yet they are
so closely packed in energy that one cannot talk of any individual state, instead
one talks of the state or level density, i.e. the number of states per unit volume
within energy E and (E + dE). The level density of allowed energy states for free
electrons in conduction band may be denoted by N(E) and for holes in valence band
by P(E). Carrier concentration for electrons and holes in a semiconductor may be
theoretically calculated using the tools of quantum statistics. Let us first calculate the
number density (concentration) of electrons in conduction band of a semiconductor.
Let us denote by n e (E) the density of electrons (number density or number of
electron per unit volume) in conduction band at energy E and (E + dE). This number
density of electrons may be written as the product of the density of energy state at
energy E and (E + dE) in conduction band, and the probability F(E) that the energy
range E and (E + dE) is occupied by electrons. Therefore,
n e (E) = N (E).F(E) (2.21a)

To calculate the number density in conduction band of electrons of all energies

n e , one must integrate Eq. (2.21a) from the lower energy limit of E C , energy at the
bottom the conduction band and E top , the energy at the top of the conduction band.
Hence,
∫Etop ∫Etop
ne = n e (E)dE = N (E)F(E)dE (2.21b)
EC EC
Exact calculation of factor N(E) for electrons is impossible as in an actual semi-

conductor crystal electrons face a periodic potential due to lattice ions. Approximate
value for N(E) assuming that electron behaves as a free particle in a box is given by
3/2
2m ∗n
N (E) = 4π E 1/2 (2.21c)
h2
Here m ∗n and h are, respectively, the effective mass of the electron and Planck’s
constant. The probability F(E) that electron occupies the state of energy E is given
by the Fermi–Dirac distribution function of quantum statistics as
1
F(E) = (2.21d)
( E−E f )
1+e
kT
In Eq. (2.21d) E f stands for Fermi energy and k for Boltzmann constant while T
is temperature in Kelvin. It is easy to verify that;
⎧
⎪ E(F) = 1 for E < E f All states with energy less than
⎪
⎨
Fermi energy are filled with electrons
At T = 0 K
⎪
⎪ E(F) = 0 for E < E All states with energy larger than
⎩ f
Fermi energy are empty
It follows from Eq. (2.21d) that for E = E f ; F(E) = 1+e 1

0 = 2 = 0.5; that
1
means the probability is 0.5 that the state with energy E f is filled with electron.
Figure 2.26i shows the variation of the function N(E) (density of states with energy
E and E + dE) with energy E. This graph basically represents Eq. (2.21c). It may
be observed in this figure that N(E) increases as the square root of energy E.
Variation of occupation probability function F(E) with energy E is shown in
Fig. 2.26ii. It may be observed that function F(E) behaves differently, for two situ-
ations T = 0 K and T > 0 K. As shown in the figure, when temperature T = 0 K,
F(E) = 1 for E < E f and F(E) = 0 for E > E f , (red curve); there is a sharp cut
in the probability at Fermi energy. For T > 0 K; there is not a sharp cut in F(E), it
has a value 1 up to some energy, then starts decreasing, becomes 0.5 at E = E f and
decreases almost exponentially after that.
Fig. 2.26 Graphical representation of N(E), F(E) and electron concentration ne at T = 0 K and T
>0K
Figure 2.26iii, iv show the number per unit volume of electrons of all energies ne
at T = 0 K and at T > 0 K, respectively. As shown in these figures, the value of ne is
given by the shaded area enclosed between the energy axis (X-axis) and curves for
functions N(E) and F(E).
E−E f
For the case where (E − E f ) > 3 kT, 1 may be neglected in comparison to e kT
in Eq. (2.21d) and F(E) may be written as

E−E f
−
F(E) = e kT
(2.21e)
Putting this value in Eq. (2.21b) one gets;
∫Etop ∫∞ ∗ 3/2
2m n −
E−E f
ne = n e (E)dE = 4π E 1/2 e kT
dE (2.21f)
h2
EC EC
The upper limit of integration in above expression is changed from E top , (energy
at the top of conduction band) to ∞ as F(E) approaches to zero exponentially for
large energies. Equation (2.21f) on integration gives
3/2
2m ∗n kT −
E C −E f
−
E C −E f
ne = 2 e kT
= NC e kT
(2.21g)
h2
where
3/2
2m ∗n kT
NC = 2 (2.21h)
h2
Similarly, number density of holes n p may be given as

3/2
2m ∗p kT −
E C −E f
−
E C −E f
np = 2 e kT
= Nv e kT
(2.21i)
h2
Here,
3/2
2m ∗p kT
NV = 2 (2.21j)
h2
(i) Calculation of Fermi energy at T ≈ 0 K

At T = 0 K, electron concentration ne may be given as
∫E f 3/2 3/2 3/2

2m ∗n 2m ∗n 2E f
ne = 4π E 1/2
.1dE = 4π
h2 h2 3
0
Or
2/3
3 h 2 2/3 0.121h 2 2/3
Ef = √ ∗
ne = ne (2.21k)
16 2π m m∗
2.4.13 Hall Effect
American physicist Edwin Herbert Hall in 1879 observed that when an electric
current is passed through a conductor that is placed in a magnetic field, a potential
proportional to the current and the magnetic field develops across the conductor in
a direction perpendicular to both the current and the magnetic field. This effect is
called the Hall Effect and the developed potential difference as Hall voltage. From
the measurements he made, Hall for the first time was able to determine the charge
of the current carriers. Even today Hall Effect is used to steady the charge transport
characteristics of metals and semiconductors.
Layout of experiment for the study of Hall Effect is shown in Fig. 2.27. A slab of
the conductor/semiconductor of length L in X direction, width w in Y-direction and
thickness t in Z-direction is taken, and a voltage source V is connected to the two
opposite faces so that an electric field ε (= V /L) in direction X is produced within
the slab. Electric field ε establishes a current I x in positive X-direction through the
slab. Current I x may be constituted by the flow of charges of only one polarity (in
case of metallic block by electrons) or it may be produced by charges of opposite
polarities (in case of semiconductor both electron and holes). However for simplicity
we assume that the current I X is due to charge carriers of only one polarity. Let q be
the charge of the carrier. The electric field ε exerts a force qε on each charge carrier
in positive X-direction and imparts an additional average velocity vX , called the drift
velocity, to each charge carrier. If n represents the concentration of charge carriers,
then the current density jx may be written as
jx = nqvx and
current Ix = nqvx .(area of cross section of the slab) = nqvx w.t

Or
Ix
vx = (2.22a)
nq.w.t
Fig. 2.27 Layout of Hall Effect experiment

If a magnetic field Bz is now applied in Z-direction, the charge carriers that consti-
tute current I x will experience a force in Y-direction. The direction of force is given by
Fleming’s left hand rule (see inset in the figure) and the magnitude by the expression
FY = q B Z vx
Force F Y will deflect charge carriers towards the top of the slab resulting in accumu-
lation of charge carriers on the inside of the top surface of the slab. Accumulation of
charges on the inner top face generates an electric field E Y in negative Y-direction.
Electric field E Y will repel charge carriers and will oppose further accumulation of
charges. Thus charge carriers will experience two opposite forces, one in positive
Y-direction F Y due to magnetic field and the other in negative Y-direction due to
electric field E Y generated by the accumulation of charges. Ultimately a state of
equilibrium will reach when two opposite forces will become equal, and no further
deflection of charge carriers will take place. In the state of equilibrium,
nq Bz vx = nq E Y or B Z vx = E y (2.22b)
Electric field E Y that is produced by the accumulation of deflected charge carriers,

results in the development of a voltage V H , called Hall voltage, across Y-direction
and may be measured experimentally. Further,
∫0 ∫w
VH = E Y dy = − E Y dy = −E Y w
w 0
Substituting the value of E Y from Eq. (2.22b) and of vx from Eq. (2.22a) in the
above expression one gets

1 B Z .Ix B Z .Ix
VH = − = −R H (2.22c)
nq t t
In expression (2.22c) V H , BZ and I X are all measurable quantities, and therefore,

in an experiment one can determine the value of RH , called Hall coefficient. RH will
have positive value if the charge of the carrier is positive and will be negative for
charges of negative polarity. For n-type semiconductor slab RH will have negative
value and for p-type semiconductor slab RH will be positive.
2.4.14 p–n Junction
When two ends of the same intrinsic wafer or intrinsic semiconductor monocrystal
are doped, one with n-type impurity of pentavalent atoms and the other by p-type
trivalent impurity, a p–n junction is formed at the boundary of the two sides. It may
be emphasised that if a p-doped crystal and another n-doped crystal are put together
touching each other, p–n junction will not be formed. For p–n junction to form it is
essential that same crystal or wafer be doped on one side by p-type impurity and on
the other side by n-type impurity then only a p–n junction is formed at the boundary
of the p- and n-type materials within the given wafer or crystal.
As is shown in Fig. 2.28, the Fermi level for isolated n-type material is shifted
upwards from the middle of the forbidden gap and it is shifted downwards for the
isolated p-type material. However, when p- and n-type materials are developed on
the same crystal, the Fermi level cannot be different on two sides because of the
continuity of crystal structure. As a result the band structure of n-side is pulled down
with respect to the band structure of the p-side to equalise the Fermi levels of the
two sides, as shown in Fig. 2.29.
The band structures of the p- and n-sides in a single crystal are shown separated
from each other in Fig. 2.29 just to indicate how the band structure of n-side is pulled
down with respect to the p-side to equalise the Fermi level on two sides. However in
reality the two band structures touch each other at the physical boundary of the p-
and n-sides.
Since the energy of free electrons in conduction band is measured up wards
from the Fermi level, electrons in conduction band on p-side are more energetic as
compared to the electrons in the conduction band on n-side. Similarly, holes on n-side
have more energy than holes on p-side.
Consider the instant when p–n junction got established on doping the two sides.
Initially both the n-side and the p-side were electrically neutral, however, at the
establishment of junction, concentration of electrons on n-side is larger than the
concentration of electrons on p-side and similarly, the concentration of holes on p-
side is larger than of holes on the n-side. Because of the concentration difference,
some electrons diffuse from n-side to p-side, and some holes diffuse from p-side to
n-side. As a result of diffusion of electrons from n-side, the n-type semiconductor
develops a positive charge; the amount of positive charge developed on n-side is
proportional to the number of electrons lost by it due to diffusion. The positive
Fig. 2.28 Band structures of isolated intrinsic, n-type and p-type semiconductors
Fig. 2.29 Energy band structures of the p- and n-sides at pn junction
charge acquired by n-side try to pull back negatively charged electrons and tries to
stop further diffusion of electrons. Thus two opposite forces; force of diffusion that
tries to transfer electrons from n-side to p-side and the force of attraction between
electrons and positively charged p-side got balanced after the diffusion of some
electrons from n-side to the p-side. This is called the state of equilibrium, in state of
equilibrium that occurs after the diffusion of some electrons from n-side to p-side,
there is no further diffusion of electrons from n-side to p-side.
As already mentioned, initially some holes, which are majority carrier on p-side,
diffuse to n-side, making p-side negatively charged. The amount of negative charge
developed on p-side is proportional to the number of holes that have diffused to
n-side. Again, at the state of equilibrium, that occurs after some holes have already
diffused to n-side, there is no further diffusion of holes.
The state of equilibrium is reached within a fraction of a second as soon as the
p–n junction is formed. After the system attains equilibrium, there is no further
diffusion of electrons from n-side and of holes from the p-side. Further, after the
establishment of equilibrium, the p-side develops a negative potential and n-side a
positive potential. The potential difference between the n-side and the p-side is called
internal potential barrier and is denoted by V B (see Fig. 2.30).
(i) Depletion layer Diffusion of electrons from the n-side leaves a sheath of uncov-
ered positive immobile donor ions on the n-side of the junction and diffusion
of holes a layer of immobile uncovered negative acceptor impurity ions on the
p-side of the junction. Thus around the junction there is a layer of positive
uncovered ion on the n-side and a layer of uncovered negative ions on the p-
side, this region which contains uncovered ions is called depletion layer. As is
Fig. 2.30 p–n junction diode with bulk p- and bulk n-sides along with depletion layer. Internal
potential barrier V B is also shown in the figure
obvious, no mobile charge carrier, electron or hole, may stay in this depletion
region, as it will be swept by the positive or negative uncovered ions. Since no
mobile charge can stay in depletion region, i.e. it is depleted of mobile charges,
hence the name depletion layer.
Depletion layer has some special properties: (i) no free mobile charge may
stay in this region, (ii) since it has no charge carriers it is like an insulator or
has very high resistance, (iii) there are equal amounts of positive and negative
charges at the two ends of the depletion layer which in itself behaves like an
insulator; therefore, depletion layer works like a parallel plate capacitor. The
capacitance of depletion layer may be changed by applying potential drop across
p–n junction, and thus it provides a capacitor whose capacitance may be varied
by varying voltage across the junction. p–n junction is also called junction diode.
It is because of the fact that p–n junction behaves like an electron tube diode.
Since the total uncovered negative charge on p-side of the depletion layer must
be equal to the total uncovered positive charge on the n-side;
n p .x p = n e X e
Here, n p , n e , X p , and X e are, respectively, the concentration of acceptor impu-

rity, concentration of donor impurity, thickness or width of depletion layer on
p-side and width of depletion layer on n-side. When the p- and n-sides are not
doped to the same concentration (n p /= n e ) the depletion layer will extend more
towards the side of lower doping concentration as shown in Fig. 2.31.
Fig. 2.31 Depletion layer

extends more on the side of
lower doping concentration
SAQ: Which part of a p–n junction has maximum resistance and why?
SAQ: There are positive ions of donor impurity atoms on the bulk n-side and
negative ions of acceptor impurity on the bulk p-side but these ions are
covered with respective charge carriers. Why do ions become uncovered
in the depletion layer of a p–n junction?
SAQ: Can you estimate the thickness of depletion layer for normal doping.
(ii) Biasing of p–n junction diode
Biasing of a device means providing required voltages to different terminals
of the device. Figure 2.32a shows the symbol used for a p–n junction diode in
electronic circuits. A junction diode has two terminals; a terminal connected to
the p-side and the other terminal connected to the n-side. A source of voltage, a
battery, may be connected between these two terminals in two different ways.
When diode terminal attached to the p-side is connected to the positive terminal
of the battery and the n-side to the negative terminal, the arrangement is called
forward bias. However, if the p-side is connected to the negative terminal of
the battery and the n-side to the positive, the arrangement is reverse bias.
(a) Forward bias
Figure 1.32b shows the forward bias arrangement. It may be recalled that
in an unbiased p–n junction at equilibrium diffusion of charge carriers
does not take place because of the internal potential barrier V B , which
restricts any transfer of charges from one side to the other. In forward
bias arrangement the battery potential V opposes or reduces the internal
potential barrier V B . Reduction of internal potential barrier results in two
events: (i) reduction in the width of the depletion layer because on p-
side holes get repelled by the external battery potential + V and covers
some of the uncovered acceptor ions in depletion layer and similarly, on
n-side electrons get pushed into depletion region by the negative external
battery potential and cover some uncovered positive donor ions. (ii) As
a result of the reduction of internal potential barrier, some of majority
Fig. 2.32 a Symbol for p–n junction diode used in electronic circuit b forward bias junction
c reverse biased p–n junction
carrier holes from p-side and some majority carrier electrons from the
n-side start moving to the other side. This movement of majority charge
carriers constitutes a forward current I f through the circuit as shown in
Fig. 2.33.
Initially, the forward current I f increases slowly till the depletion layer
vanishes completely at battery potential V d when forward current suddenly
rises almost exponentially. Potential V d is called knee potential (or on
potential) and in a way equal to the internal potential barrier V B . For Silicon
Fig. 2.33 Characteristics of

a p–n junction diode
p–n junction diode the knee potential has a value of 0.7 V and for Ge based
p-n diode it is 0.3 V. On further increasing the forward bias voltage beyond
V d , the voltage across the junction does not increase but forward current of
larger value flows through the forward biased circuit. A p–n junction diode
in forward bias above on-voltage V d behaves as a battery of 0.7 V in case
of Silicon-based diode and a battery of 0.3 V in case of Germanium-based
diode. When forward bias voltage is increased beyond V d , the depletion
layer disappears and large number of majority carriers diffuse from both
sides to the opposite side. Therefore, forward current if is essentially due
to the diffusion of majority carriers and since the concentration of majority
carriers is quite high (≈ 107 charge carrier cm−3 ) forward current I f of few
milli amperes flows through the circuit.
(b) Reverse bias
Circuit diagram for reverse bias arrangement is shown in Fig. 2.32c. In
reverse bias arrangement the external battery potential add up with the
internal potential barrier V B . This results in the increase of the width of
depletion layer. With enhanced barrier at junction (V + V B ) the majority
carrier on the two sides do not cross the depletion layer of enhanced width.
However, minority carriers, electron on the p-side and holes on n-side are
pushed by the total barrier potential (V B + V ) across the depletion layer
constituting reverse current I r . This flow of minority carriers from one side
to the other is not due to diffusion instead it is due to the large potential
difference across the depletion layer. Minority carrier current in reverse bias
is drift current and is only of the order of few microamps. When reverse
bias voltage is increased beyond the breakdown voltage V b (see Fig. 2.33),
suddenly a large reverse current starts flowing in the circuit. This large
current flows because of the breakdown of the crystal structure. Because
of the large electric field inside the semiconductor crystal (established by
large reverse bias voltage V b ), the atoms in crystal structure break down
releasing large number of electrons.
Graphs showing the variation of forward and reverse currents as a function
of applied voltage are called p–n junction diode characteristics and are
shown in Fig. 2.33 both for the forward and the reverse bias arrangements.
SAQ: Forward current across a p–n junction is generally in mA while the
reverse current is in μA. What is the reason for this difference in
the magnitudes of the two currents?
SAQ: It is known that reverse saturation current I r is very sensitive to
the ambient temperature; for every 100 C rise of temperature it
gets doubled. However the forward current I f is not so sensitive to
temperature. Can you assign a reason for this difference?
Semiconductor materials are the backbone of electronic industry. These mate-
rials are used in fabricating solid state electronic devices that are extensively
used in modern analogue and digital electronics. p–n junctions developed in
special conditions of doping at more than one place in a monocrystal give rise
to bipolar junction transistors and field effect transistors.
2.4.15 Some Formulations
Some important formulae that are applicable in case of semiconductors are given
here without their derivations, which are beyond the scope of the present text. These
formulae may be used to solve numerical problems.
(1) If ne and np, respectively, denote the concentrations of free electrons and holes in
doped semiconductor at temperature T and ni the concentration of free electrons
or holes in the intrinsic semiconductor at same temperature T, then,
n e .n p = n i2 (2.23)
(2) For a non-degenerate semiconductor doped with shallow dopant at room

temperature,
ne ∼
= N D and n p ∼
= NA (2.24)
Here, ne and np are, respectively, electron and hole concentrations, while N D

and N A are concentrations of donor and acceptor atoms, respectively.
(3) Theoretical values of carrier densities: theoretical value for the concentration of
free electrons ne and holes np in a semiconductor at temperature T, calculated
using quantum statistics, is given as
23
−
E C −E F
2π m ∗n kT
n e = NC e kT
; NC ≡ 2 (2.25)
h2
And
23
−
E F −E V 2π m ∗p kT
n p = NV e kT
; NV ≡ 2 (2.26)
h2
Here, E C is the energy at the bottom of the conduction band, E V the energy at the
top of valence band, k Boltzmann constant, m ∗n , m ∗p , respectively, the effective
masses of electron and hole and T temperature in Kelvin.
(4) Positioning of Fermi level
Fermi energy at T ≈ 0 K is given by Eq. (2.21k) as
2/3
3 h 2 2/3 0.121h 2 2/3
Ef = √ ∗
ne = ne
16 2π m m∗
Exact calculations show that Fermi level for an intrinsic semiconductor is

positioned according to the following relation;
∗
(E C − E V ) 3kT mp
E Fint = + ln (2.27)
2 4 m ∗n
Here, E C and E V are, respectively, the energies at the bottom of conduction band
and the top of valence band. m ∗p and m ∗n are, respectively, the effective mass of hole
and electron. Second term on right that depends on temperature is negligible at room
temperature and varies slowly with temperature, therefore, it is often neglected.
Hence Fermi level for intrinsic semiconductor is taken at the middle of the forbidden
energy gap at all temperatures.
n-type
However, for n-type semiconductor, it may be shown that Fermi level E F is
given by
n-type NC
EF = E C − kT ln (2.28)
ND
Here E C is the energy at the bottom of conduction band, k Boltzmann constant, T is

3
2π m ∗ kT 2
temperature in Kelvin, NC ≡ 2 n
h2
and N D is concentration of donor impurity
atoms. Equation (2.28) tells that Fermi level in n-type semiconductor shifts towards
the bottom of conduction band with the increase in donor impurity concentration.
p-type
Similarly, for p-type material the position of Fermi level E F is given by
n-type NV
EF = E V + kT ln (2.29)
NA
Here E V is the energy at the top of valence band, N A concentration of acceptor

3
2πm ∗ kT 2
impurity atoms and N V ≡ 2 p
h2
. Equation (2.29) tells that the Fermi level for
p-tpye material shifts towards the top of valence band with the increase in acceptor
impurity concentration.
Values of Some Important Constants

Energy 1.6 × 10–19 J = 1.0 eV.
Thermal energy corresponding to room temperature (300 K) = kT = 0.025 eV.
Unit of charge e = 1.6 × 10–19 C.
Mass of electron me = 9.1091 × 10–31 kg.
Planck’s constant h = 6.62608 × 10–34 J s = 4.1357 × 10–15 eV s.
Velocity of light c = 2.9979 × 108 m s−1 = 3.0 × 108 m/s.
Boltzmann constant k = 1.3807 × 10–23 J K−1 = 8.6173 × 10–5 eV K−1 .
Solved Example SE2.1 A semiconductor absorbs light photons of wavelength

shorter than 1 μm and is transparent to photons of wavelengths larger than it. What
is the magnitude of forbidden energy gap of the material?
Solution: Light photons that are absorbed by the semiconductor will transfer valence
band electrons to conduction band. It is obvious that the forbidden energy gap will
be equal to the energy of the photon. Now photon energy E is given by

hc 4.1357 × 10−15 (eV s) × 3.0 × 10−8 m s−1
E = hν = = = 1.24 eV.
λ 1 × 10−6 (m)
Therefore the forbidden energy gap E g = 1.24 eV.

Solved Example SE2.2 The forbidden energy gap for a semiconductor crystal at
temperature 300 K is 1.50 eV. Take the effective mass of electron and hole, respec-
tively, as 0.1 me and 0.5 me and calculate the energy shift of Fermi level from the
middle of the forbidden energy gap.
Solution: We use Eq. (2.27)
∗
E Fint = + ln
2 4 m ∗n
Energy shift of Fermi level from middle of forbidden energy gap is

∗
E Fint − = ln
2 4 m ∗n
3 × 8.61 × 10−5 × 300 0.5m e
= ln = 19.37 × 10−3 × 1.609
4 0.1m e
= 0.031 eV
Solved Example SE2.3 The intrinsic concentration of charge carriers in a semicon-

ductor is 1 × 1018 m−3 . Calculate the conductivity of the semiconductor given that
electron and hole mobilities are, respectively, 0.40 and 0.15 m2 V−1 s−1 .
Solution: Conductivity
σ = n i .e.[μe + μh ] = 1 × 1018 × 1.6 × 10−19 [0.40 + 0.15]

= 0.088 Ω−1 m−1 .
Solved Example SE2.4 Forbidden energy gap for a compound semiconductor is

1.4 eV at 300 K. The Fermi level for the doped semiconductor is shifted towards the
valence band by 0.20 eV. What is the type of doping and what is the majority carrier
concentration? Given that effective mass of electron is 0.06 me and of hole 0.5 me .
Also calculate the concentration of minority carrier in doped semiconductor and the
concentration of charge carriers in intrinsic semiconductor.
Solution: It is given that after doping the Fermi level shifts towards the valence band.
It means that the doping is done with acceptor impurity and that the material has
become p-type after doping.
The band structures of the semiconductor before doping (intrinsic material) and
after doping (p-type material) are shown in Fig. 2.34. As indicated in the figure after
doping (E F − E V ) = 0.5 eV and (E C − E F ) = 0.9 eV.
To calculate concentration of majority carrier holes we use Eq. (2.26) given below;
23
−
E F −E V 2π m ∗p kT
n p = NV e kT
; NV ≡ 2
h2
Let us first calculate the value of N V
23 23
2π m ∗p kT 2π × 0.5 × 9.1 × 10−31 × 1.38 × 10−23 × 300
NV ≡ 2 =2 2
h2 6.63 × 10−34
= 8.836 × 1024 .
Next we calculate

E F −E V 0.5×1.6×10−19
− −
n p = NV e kT
= 8.836 × 1024 e 1.38×10−23 ×300 .
Or
n p = 8.836 × 1024 × e−19.32 = 8.836 × 1024 × 4.068 × 10−9 = 3.59 × 1016 m−3 .
Fig. 2.34 Band structure before and after doping

Next we calculate the concentration of minority carrier electrons using Eq. (2.25)
given below,
23
−
E C −E F
2π m ∗n kT
n e = NC e kT
; NC ≡ 2 .
h2
Now,
23 3/2
2π m ∗n kT 2π × 0.06 × 9.1 × 10−31 × 1.38 × 10−23 × 300
NC ≡ 2 =2 2
h2 6.63 × 10−34
= 2.70 × 1021 .
Or
E C −E F
−19

− − 1.38×10
0.9×1.6×10
n e = NC e kT
= 2.70 × 1021 × e −23 ×300
= 2.70 × 1021 × e−34.78
Or
n e = 2.70 × 1021 × 7.8568 × 10−16 = 2.12 × 106 m−3 .
Concentration of charge carriers in intrinsic material ni is given as

2 √
n i = n e n p = 2.12 × 106 × 3.59 × 1016 = 2.76 × 1011 m−3 .
Therefore,
(i) majority carrier concentration n p = 3.59 × 1016 m−3
(ii) Minority carrier concentration n e = 2.12 × 106 m−3 and
(iii) Charge carrier concentration n i = n ip = n ie = 2.76 × 1011 m−3 .
Solved Example SE2.5 Given that number density of free electrons in gold at very
low temperature ≈ 0 K is 6.0 × 1022 cm−3 , calculate the Fermi energy for gold. Take
the effective mass of electron to be equal to its mass 9.1 × 10–31 kg.
Solution: In the given problem number density of electrons is given in CGS units
while the electron mass is in MKS units. Let us convert electron number density also
in MKS units; given quantities are;
n e = 6.0 × 1022 cm−3 = 6.0 × 1028 m−3 , the effective electron mass m ∗e =
9.1 × 10−31 kg.
We use the formula given by Eq. (2.21k) for Fermi energy E F at T = 0 K.
2 2/3
E F = 0.121×h
m∗
n e which on substituting the values gives,
e
2
0.121 × 6.60 × 10−34 2/3
EF = −31
6.0 × 1028 J
9.1 × 10
2.5 Conductors 113
8.87 × 10−19
= 8.87 × 10−19 J = eV = 5.548 eV.
1.6 × 10−19
2.5 Conductors
Conductors are solids that are characterised by metallic bonding, having either over-
lapping conduction and valence bands or with negligible forbidden energy gap.
Metals and their alloys are mostly conductors. Their specific resistivity lies in the
range of (1–100) × 10–8 for metals and (1–100) × 10–6 Ω m for most of the alloys.
Overlapping of conduction and valence bands is the outcome of the high degree of
overlap in outer electron orbital’s of individual atoms in some crystals. As a result the
conduction and valence bands become so broad that they overlap. In such materials
the valence electrons are far away from the corresponding nucleus of the atom and
are very loosely bound with its parent nucleus. Also in their crystalline structure the
relative separation of atoms is large so that the forbidden energy gap is either zero
or very small.
Figure 2.35 shows the band structure of a conductor (a) at 0 K and (b) at T > 0 K.
At absolute zero all valence electrons are bound and are not available for conduction
of current. However, with the rise of ambient temperature more valence electrons
become delocalised and at room temperature in most of conductors all valence elec-
trons become delocalised or free and are available for conduction of current. It is
reasonable to assume that at room temperature all valence electrons of all atoms in the
given specimen of conductor are delocalised and are available as free charge carriers.
Obviously, current may flow only when some voltage is applied to the conductor that
establishes an electric field. In absence of any electric field a piece of conductor at
room temperature has large number of free electrons that move in random directions
with thermal velocity which is of the order of 105 m/s. The number density of free
electrons in Silver at temperature 300 K is of the order of 5.8 × 1028 m−3 . These
randomly moving electrons undergo frequent collisions with crystal lattice and are
also trapped at sites of unionised impurity atoms. At each collision the velocity and
direction of motion of electron get changed. When randomly moving free electrons
are subjected to an electric field by applying an external voltage, a drift velocity gets
superimposed on each electron in a direction opposite to the direction of the elec-
tric field. This results in the flow of current through the conductor. Opposition to the
smooth flow of current is generated by frequent lattice-electron collisions. Larger the
rate of collision more will be the opposition to the flow of current. Therefore, resis-
tance or resistivity of conductors is essentially the result of electron–lattice collisions.
With the rise of temperature, the electron density in a conductor does not increase
because all valence electrons are already delocalised at room temperature, however,
the electron–lattice collision rate increases with temperature and hence the resistivity
of conductors increases with temperature. That is why the temperature coefficient of
conductors has a positive value. On the other hand, in a semiconductor, all valence
Fig. 2.35 Band structure of a conductor a at absolute zero temperature and b at a temperature
higher than absolute zero
electrons are not free at room temperature, and the free electron density rapidly rises
with temperature due to covalent bond breaking. Though electron–lattice collisions
in semiconductors also rise with temperature, but the rate of increase of free charge
carriers (electrons and holes) with temperature is much larger than the increase of
collisions; therefore, the resistivity of semiconductors decreases with temperature.
That is why semiconductors have negative value of temperature coefficient.
2.5.1 Semimetals and Half Metals
(i) Semimetal or metalloid

Semimetals are materials for which the bottom of the conduction band either just
touches the top of the filled valence band or there is very little overlap between
the two bands. As a result, there is a very small value for the density of elec-
tron states near Fermi level in semimetals, as compared to metals where there
is considerable overlap of bands (see Fig. 2.36). With few energy states avail-
able around Fermi level, semimetals behave as semiconductors with negligible
forbidden energy gap. Some semimetals, also called metalloids, are arsenic
(As), Antimony (Sb) and tellurium (Tl). Some elements like carbon are found
in two allotropic forms, one of which behaves like a metal (graphite) while the
other (diamond) like a non-metal, such elements are also included in the list of
metalloids
(ii) Half metal
It is often said that a given electron level can accommodate at the most two
electrons, one with spin up and the other with spin down. However, on close
examination it is found that each electron level is made up of two very closely
placed electron states, one for spin up electron and the other for spin down
2.5 Conductors 115
Fig. 2.36 Energy band structure of a metal b semimetal
electron. As such one can build two separate state energy diagrams one for spin
up electrons and the other for spin down electrons with their own valence and
conduction bands.
In some crystals where atoms are bound by metallic bonding, it so happen that
valence band for electrons of one specific spin orientation is partially filled
while there is a forbidden energy gap for electrons of other spin orientation.
As a result when external voltage is applied, electrons with that spin orienta-
tion for which there are vacant states in valence band contribute to the flow of
current. Electrons with opposite spin orientation do not contribute to current
flow because of forbidden energy gap. Since only about half of the total elec-
trons contribute to current flow, the material is termed as half metal. Examples
are chromium oxide and lanthanum-strontium-magnetite that are half metals
and are also ferromagnetic. Though all half metals are ferromagnetic but all
ferromagnetic materials are not half metals. Energy band structure of a typical
half metal is shown in Fig. 2.37.
Solved Example SE2.6 Density of trivalent Aluminium metal is 2.7 g cm−3 , and its
molecular mass is 27 g/mol; assuming that at room temperature all valence electrons
are non-localised (or free), calculate the number density of free electrons in the metal.
Solution: It is known that a gramme mole of an element contains 6.022 × 1023 atoms
(Avogadro’s number) of the element. Therefore,
27 g of Al will contain 6.022 × 1023 atoms of Aluminium. (2.30)
Also, the density D of Al is given as D = 2.78 g cm−3 . But density is equal to mass/
volume.
Fig. 2.37 Energy band structures for a half metal a for spin down electrons b for spin up electrons
We calculate the volume V of 27 g of Aluminium using its density as
M 27
V = = cm3 = 9.71 cm3 (2.31)
D 2.78
It follows from Eqs. (2.30) and (2.31) that 27 g of Aluminium has 6.022 × 1023 atoms
that occupy a volume of 9.71 cm3 .
6.022×1023
Therefore, the number of atoms in 1 cm3 = 9.71
= 0.62 × 1023 .
Since the valency of Aluminium is 3, and all valence electrons are free at room temper-
ature, therefore each atom will contribute three free electrons at room temperature.
Hence the number density of free electrons in Aluminium at room temperature n per
cc is
n = 3 × 0.62 × 1023 per cm3 = 1.86 × 1023 per cm3 = 1.86 × 1029 per m3 .
2.6 Superconductor
Resistivity of conductors, particularly of metals is very low of the order of 10–8

Ω-m, and it originates essentially from the free electron–lattice or electron–phonon
interactions. It is because of the resistivity that energy is lost when current is passed
through a conductor. As expected, resistivity of conductors decreases further with the
decrease of temperature. In some materials, generally at very low temperature near
about absolute zero, it is found that their resistivity simply vanishes. The state of the
2.6 Superconductor 117
material in which the resistivity of the material becomes zero is called the supercon-
ducting state and the property as superconductivity. The characteristic temperature
below which resistivity becomes zero is called the critical or transition tempera-
ture and is denoted by TC . It is obvious that no energy loss occurs when current is
established through a superconductor, no matter for how long current flows through
it.
2.6.1 Background
Research group of Dutch physicist Heike Kamerlingh Onnes, in 1911 found that the
resistivity of a mercury column becomes zero when the temperature of the spec-
imen was reduced below 4.15 K (see Fig. 2.38). Complete disappearance of electric
resistance in some other metals and solids below a certain characteristic very low
temperature was observed in some other materials also.
Onnes and his students were studying the electrical behaviour of wires of different
materials and found that the resistance of a mercury wire took a precipitous drop
when temperature reached to about 4.15 K. The drop in resistance was enormous,
the resistance of the wire dropped at least by a factor of one thousand, so much
so that exact measurement of the resistance became impossible (see Fig. 2.38). In
order to further investigate the phenomenon, Onnes’ group setup a current through
the mercury wire in the form of a ring by connecting at two points of the wire a
voltage source for an instant and then removing the voltage source. To their surprise,
they observed that current kept flowing through the mercury wire ring without any
reduction in its magnitude, so long as the temperature of the wire was kept below
4.15 K. The observed perpetual flow of current was only possible if flow of current
does not encounter any opposition or resistance. As is known opposition to the flow
of current in normal situation arises essentially from electron–lattice collisions and
electron trapping at impurity sites. Disappearance of resistance in case of mercury
Fig. 2.38 Superconducting

transition at critical
temperature T c = 4.15 K
Table 2.3 Superconducting

Metals Transition temperature (K)
transition temperatures for
some metals Lead 7.19
Mercury 4.15
Tin 3.72
Indium 3.41
Aluminium 1.20
Zinc 0.88
wire at temperature below 4.15 K means that electron–lattice collisions have either
suddenly cease to happen below the critical temperature, the temperature below
which mercury wire exhibits superconductivity or at least lattice vibrations are not
opposing the flow of current. Transition or critical temperature for superconductivity
transition for some metals is listed in Table 2.3.
Many well-known scientists including Nobel Lauriat John Bardeen tried to explain
and give a theoretical background for superconductivity but they did not succeed.
The reason why no theoretical explanation of the process could be given at that time
was that the process of superconductivity is a typical quantum phenomenon, and
quantum physics was not in place till 1920 or so.
(i) Meissner effect In the meantime experimental studies on superconductivity
continued and in 1933, two scientists Walter Meissner and Robert Ochsen-
feld discovered another very interesting property of superconductivity; they
found that any material in superconducting state repels the lines of external
magnetic field (Bex ) so long as the applied magnetic field is below the crit-
ical value denoted by BC ex . It means that for external magnetic fields Bex <
Bc ex a superconductor behaves as a perfect diamagnetic material. If a magnet
is brought near to a superconducting material, the superconductor does not
allow magnetic lines of force to penetrate through it, rather it repels them. The
effect is called Meissner effect.
Meissner effect is a typical example that also shows that superconductors are
not just perfect conductors. A perfect conductor may be defined as a conductor
which has a pure crystalline structure without any impurity or missing atom
sites and has small value of resistivity. However, a superconductor is different
from a perfect conductor as it behaves differently than a perfect conductor
when a magnetic field is first applied and then switched off.
Figure 2.39 shows a perfect or ideal conductor and a superconductor, initially
the temperature of both the specimen is above critical temperature T C and both
of them allows the passage of magnetic lines through them. With magnetic
field (Bex < BC ex ) on, if the temperature of both specimens is reduced below
critical temperature T C , the ideal conductor will allow magnetic lines to pass
through it, as they were before the reduction of temperature. It is because
magnetic properties, like susceptibility, of a perfect conductor do not change
with temperature. However, in the case of superconducting specimen, reduc-

tion of temperature below T C transforms it into superconducting state. The
magnetic susceptibility of superconducting state is different than that of the
non-superconducting state. As a result, at the instant when temperature goes
below T c , magnetic flux linked with the superconducting volume changes.
This change in the magnetic flux induces surface currents at the outer skin
of the superconducting volume such that a magnetic field exactly equal in
magnitude but opposite in direction to the previously existing magnetic field is
generated. The previously existing magnetic field gets cancelled by the induced
magnetic field in the interior volume of the superconducting specimen. Hence
no magnetic field stays in the superconducting volume. The induced magnetic
field cancels the existing magnetic field so that the interior of the supercon-
ducting volume becomes free of all magnetic fields; no magnetic field and flux
remain linked with the superconducting volume.
Let us now discuss what happens when external magnetic field is switched off
keeping the temperature T < T c .
At the instant when external magnetic field is switched off, the flux linked
with the ideal conductor goes to zero from the initial value Bex . The changing
Fig. 2.39 Behaviour of perfect or ideal conductor and superconductor with respect to magnetization
and demagnetization
magnetic flux induces surface current at the outer skin of the ideal conductor,
which in turn establishes a magnetic field in the interior volume of the perfect
conductor specimen. It may, however, be mentioned that the induced surface
currents will be short lived as the resistivity of the conductor will dissipate
energy, and currents will die out.
In case of the superconducting specimen, no magnetic flux is lined with the
specimen volume (as there is no magnetic field inside superconducting volume)
hence at the instant when Bex is switched off no change in magnetic flux will
take place. As such no induced currents will be generated. The interior and
exterior of the superconducting volume will contain no magnetic fields after
the external magnetic field is switched off.
SAQ: How can one explain the total absence of magnetic field in the interior
of a superconducting volume when some external magnetic field is
applied to the superconductor?
SAQ: When an external magnetic field is switched off from a normal
conductor, the conductor retains magnetic field in its interior and
around. How one can explain this retention of magnetic field?
(ii) Magnetic field trapped in a superconducting ring
Figure 2.40i shows a ring shaped superconductor specimen placed in an
external magnetic field Bex at temperature T > T c . Science temperature is
above T C , the specimen ring behaves as a normal material and magnetic lines
penetrate through the opening of the ring. The magnetic flux ϕ linked with the
opening of the ring is given by
ϕ = Bex .(Area of the ring opening)
In the next step temperature T is reduced below the critical temperature T C

and the magnetic field Bex is switched off. With the decrease of temperature T
below T c , the specimen ring becomes superconducting while switching off of
magnetic field (Bex = 0) will induce an electric field E in the ring according
to Faraday’s law.
∮
Edl = − dϕ dt
, here E is the electric field along the closed loop of supercon-
ducting ring and ϕ is the magnetic flux through the opening of the ring. But
no electric field can exist in a superconductor, hence,
∮
Edl = − dϕ dt
= 0, that means the flux linked with the ring before it becomes
superconducting will continue to remain linked with the opening of the ring
even after switching off of the magnetic field as shown in Fig. 2.40ii.
SAQ: Why electric field cannot sustain in a superconductor?
(iii) Superconductor type-I and type-II
Fig. 2.40 Magnetic field trapped in the opening of a superconducting ring
It is known that not only metals but some other materials below their transition
temperature Tc become superconductor. Further, if an external magnetic field
Bex is applied across a superconductor specimen, when it is below its transition
temperature, magnetic field Bin inside the superconductor stays zero. This is,
however, true only when the magnitude of the externally applied magnetic field
is below a certain value Bc ex . If the magnitude of externally applied magnetic
field Bex is increased beyond the critical value Bc ex then the superconductor
may respond in two different ways, depending on its type. In the case of
type-I superconductor on increasing the strength of external magnetic field
Bex beyond Bc ex , the superconductivity of the specimen just vanish, though its
temperature is still below Tc . It behaves as an ordinary conducting material
and magnetic field penetrates in the interior of the specimen. This is shown in
Fig. 2.41a where a dotted vertical line at Bc ex divides the figure in two parts;
where the specimen remains a superconductor and the part where superconduc-
tivity is totally lost in spite of its temperature being below critical temperature
T c . Figure 2.41a shows the typical behaviour of a type-I superconductor. Most
metals show type-I superconductivity.
In case of type-II superconductors, there are two values of external critical
magnetic fields Bc1 ex and Bc2 ex such that between these two magnetic field
values the specimen remains partially superconducting as shown in Fig. 2.41b.
For external magnetic fields greater than Bc2 ex , the type-II specimen also
becomes non-superconductor, though its temperature is still below its critical
temperature.
Fig. 2.41 a Type-I semiconductor. b Type-II semiconductor
During the phase of partial superconductivity (between Bc1 ex and Bc2 ex ),

cylindrical tubular regains in type-II specimen become ordinary or non-
superconductor as shown in Fig. 2.42. External magnetic field penetrates
only in these cylindrical volumes which are almost uniformly distributed over
the specimen volume. Initially, just beyond Bc1 ex , the non-superconductor
(normal) cylindrical volumes are very thin; however, their radius increases
as external magnetic field is increased from Bc1 ex, and finally the non-
superconducting material fills the total volume of the specimen at Bc2 ex and
beyond. The origin of these normal-filament like structures is beyond the scope
of the present text.
Table 2.4 give a list of some compounds that are type-II superconductors.
In most of type-II superconductors the value of the external field Bc2 ex is quite
high therefore, they may withstand high magnetic fields without losing super-
conductivity altogether. Wires made of type-II superconductor are frequently
used to build powerful electromagnets. For example, the upper critical external
magnetic field Bc2 ex for wires of niobium-tin (Nb3 Sn) is as high as 24 T. These
wires are used in making strong electro magnets for use in MRI and other
imaging machines. The advantage of superconducting electromagnets is that
current only has to be applied once to the wire, which are then formed into
a closed loop and allow the current and the magnetic field to persist indefi-
nitely as long as the temperature is kept below the critical temperature. Once
current is established in superconducting magnet, the external power supply
may be switched off. On the other hand the strongest magnetic field that may
be produced using a permanent magnet may be only of the order of few Tesla.
SAQ: What are type-II superconductors?
Fig. 2.42 Normal or non-superconducting cylindrical volumes in type-II superconductor. External

magnetic field penetrates through these filaments like normal volumes. The radius of these filaments
increases with the increase of external magnetic field
Table 2.4 Compounds that are superconductor

Compound T C (K) Compound T C (K) Compound T C (K)
Nb3 Sn 18.1 PbMo6 S8 15.0 YPd2 B2 C 23.0
Nb3 Ge 23.2 HoN12 B2 C 7.5 UPt3 0.5
Cs3 C60 19.0 UPd2 Al3 2.0
MgB2 39.0 (ET)2 Cu [Ni(CN)2 ]Br 11.5 (TMTSF)2 ClO4 1.2
(iv) Stable levitation

In order to demonstrate the Meissner effect, a high-temperature superconductor
like YBa2 Cu3 O7 is taken and cooled below its critical temperature (93 K). A
small and strong permanent magnet is then placed on top of the superconductor
to show the repulsion of the small magnet by the superconductor. Repulsion
results in the levitation, or hanging/floating of the small magnet above the
superconductor. Repulsion originates from the mirror reflection of perma-
nent magnets lines of force by the diamagnetic superconductor that forms an
inverted mirror image of the permanent magnet. The situation is comparable
to placing one magnet over another identical magnet to achieve the floating of
the one over the other. In principle levitation of a small permanent magnet over
the superconductor should be possible only if the size of the superconductor
is much larger than the size of the permanent magnet. In case the two are of
the same size, distortion of magnetic lines at the rim of the magnet and in its
mirror image at superconductor will make levitation unstable and the magnet
will topple down, instead of levitating. This is exactly what happens in case
of two identical magnets, stable levitation is never achieved. However, stable

levitation in case of a small permanent magnet and a superconductor of only
slightly larger in size than the permanent magnet does take place. How can
this be explained? Not only the levitation is stable but a small nudge causes
the magnet to spring back to its original position as if some unseen springs are
holding the magnet at its position of levitation.
The answer is hidden in the properties of type-II superconductor and the trick
done by the one who carries out the levitation experiment. A keen observation
of the levitation experiment will reveal that levitation of small permanent
magnet is not stable if a freshly cooled superconductor is taken and the magnet
is made to float over it; the magnet falls off in a short time. However, to
make levitation stable, the experimenter pushes the permanent magnet towards
the superconductor and then the levitation over a superconductor of small
dimensions becomes stable. So the trick to achieve stable levitation is to make
the small magnet thrust towards the superconductor.
The superconductor taken to demonstrate levitation is type-II supercon-
ductor, pushing or thrusting the permanent magnet towards the superconductor
increases the strength of the external magnetic field applied to the super-
conductor beyond the first critical value Bc1 ex . The type-II superconductor
goes into the region of partial superconductivity. Large number of filament
like normal (non-superconducting) regions develops in the superconductor.
Magnetic lines of force originating from the permanent magnet penetrate in
these filament shaped non-superconducting regions as these regions are not
diamagnetic, they are normal regions. Penetrating magnetic force lines provide
a cluster of strings that holds and brings back the floating magnet if displaced
from its position of levitation.
Levitation by small type-II semiconductors is robust and very stable. Crucial
role in it is played by the filament like normal regions that develop in type-II
superconductors in the region of partial superconductivity. These filaments like
structures are called by many different names; flux lines, fluxoids, vortices,
fluxons, etc. In a pure monocrystal of a type-II superconductor it can be shown
that these vertices are flexible, i.e. they may bend and may not be straight.
However, in a superconducting crystal that has impurities and internal struc-
tural boundaries, vortices become straight and rigid. This is called flux or
vortices pinning. Flux pinning also play important role in stable levitation.
SAQ: In your opinion which type of superconductor is more useful and why?
SAQ: What is pinning? Discuss the role played by pinning in stable
Lavitation.
(v) High Tc superconductors
Since the discovery of superconductivity it has been the desire of all scientists to
develop materials that may show superconductive behaviour at room temper-
ature. Room temperature superconductors would have revolutionise almost
every aspect of modern technology, particularly, power generation, its trans-

portation, fabrication of high magnetic field electromagnets, etc. However, till
date room temperature superconductors have not been fabricated or discov-
ered. Two scientists, Georg Bendnorz and K. Alex Muller, working at IBM
lab, discovered in 1986 a class of materials that showed superconductivity at
liquid nitrogen (LN2 ) temperature. These materials called high-temperature
or high T c superconductors are frequently used to demonstrate supercon-
ducting behaviour on bench-top using cooling by liquid nitrogen, which is
easily available. The LN2 temperature is around 77 K, much higher than the
critical temperature required for older superconducting materials. Table 2.5
lists some high Tc superconductors. Bendnorz and Muller got Nobel Prize for
their discovery in 1987, within one year of their discovery. It was the fasted
award of Nobel Prize ever. In most bench-top demonstration experiments one
uses Yttrium-Barium-copper oxide (YBa2 Cu3 O7 ) with liquid nitrogen cooling.
(vi) Isotope effect The single most important experimental observation that paved
way for developing an explanation of superconductivity was the ‘isotope effect’
discovered in 1950, simultaneously at two laboratories; by Reynolds, Serin,
Wright and Nesbitt at Rutgers University and by Maxwell working at the
National Bureau of Standards. These experimentalists accurately measured
the critical temperatures for superconductive transition in mercury samples of
different isotopic mass distributions. It is known that mercury has seven stable
isotopes with mass numbers; 196, 198, 199, 200, 201, 202 and 204. Out of these
seven, the two most abundant isotopes are 200 Hg (23,1%) and 202 Hg (29.7%).
It is possible to change the relative percentage of different isotopes in different
samples, and these experimentalists measured the critical temperature for these
different mercury samples that have different values of average isotopic mass.
They found that the critical temperature was inversely proportional to the
square root of the average isotopic mass of mercury sample, see Fig. 2.43.
Similar experiments were later carried out with other materials, and the isotope
effect was observed in each case. This showed that superconductivity has
something to do with the nuclei of the atoms of the material.
Inverse dependence of critical temperature on the square root of average
mass was something that has been observed elsewhere also. For example, when
a mass attached with a spring is given a push, the spring starts vibrating, and
the frequency of vibration is found to be inversely proportional to the square
root of the mass. In a crystalline solid it may be assumed that different atoms
are connected with each other by some sort of springs and giving a push to any
atom will set the complete atomic layer, the lattice in vibratory motion. Isotope
Table 2.5 List of some high

Compound T C (K) Compound T C (K)
T c superconductors
HgBa2 Ca2 Cu3 O8+x 135 Ti2 Ba2 Ca2 Cu3 O10+x 125
Bi2 Sr2 Ca2 Cu3 O10+x 107 YBa2 Cu3 O6+x 93
Fig. 2.43 Isotope effect for

Mercury
effect indicates that current flow in superconductor is not only an electronic

phenomenon but it very much depends on lattice vibrations.
The isotopic effect showed that although the electrical conductivity was
known to arise because of the motion of free electrons and the resistivity essen-
tially due to collisions between electron and lattice vibrations, but lattice vibra-
tions below critical temperature somehow help in smooth flow of electrons
setting perpetual current in superconductors.
SAQ: What is the significance of the isotope effect of superconductor?
(vii) Cooper pair
American physicist Leon Cooper in 1956 described a process of binding of two
electrons (or any other Fermions) in crystalline solids at very low temperature.
This special type of binding between two electrons that are quite far apart
develops on account of the interaction between the lattice vibration quanta
phonon and electrons.
It may sound strange as to how an attractive force may develop between two
negatively charged electrons that should repel each other, but this happens
because of the distortion produced in electric field of crystal lattice on account
of its vibratory motion. It is known that in metals there is large number of
free electrons since almost all atoms of the solid are ionised. The positive ions
of constituent atoms are arranged in a regular pattern in 3-dimensions. A two
dimension plane containing positive ions is termed as crystal lattice. Ions of
the lattice are held at their place by electrostatic forces of mutual repulsion
and ion–electron cloud attraction. Each lattice ion is in a sort of dynamic
equilibrium under the two opposite forces. These dynamic forces between
lattice ions may be compared to small springs which hold ions in the lattice. A
little jerk or push to an ion of the lattice may make the whole lattice to undergo
vibratory motion. This vibratory lattice motion is quantized, and the quanta of
vibratory motion of the lattice are called phonon.
Figure 2.44 shows two electrons numbered 1 and 2 moving in opposite direc-
tions. Negatively charged electrons attract positive ions of the lattice towards
them, distorting the crystal lattice and setting it in vibratory motion. The density
of positive charge in regions around the lattice distortion increases beyond its
normal value and becomes centres of attraction between electrons and the
distortion. Nearby electrons feel force of attraction by the region of increased
charge density but this force of attraction is over powered by the force of
mutual repulsion between nearby electrons. However, electrons far away get
bound with each other as the force of mutual repulsion between distant elec-
trons is very small and force of attraction due to charge distortion over rides the
repulsion. In this way Cooper electron pairs are formed. This is the reason
why the distance between the two electrons of Cooper pair may range from 50
to 100 nm or more. This is in comparison with the lattice separation; distance
between two neighbouring ions of the lattice, of 0.1–0.4 nm. It is important to
note that only those electrons that are far apart may form Cooper pairs.
Figure 2.44 shows two electrons numbered 1 and 2 moving in opposite direc-
tions. Negatively charged electrons attract positive ions of the lattice towards
them, distorting the crystal lattice and setting it in vibratory motion. The density
of positive charge in regions around the lattice distortion increases beyond its
normal value and becomes centres of attraction between electrons and the
Fig. 2.44 Moving electrons produce distortion in ion lattice and set it in vibratory motion. Distortion
in lattice increases density of positive charge in small regions which attract electrons that are far
away and create Cooper pairs
Fig. 2.45 Phonon

interaction between electrons
of Cooper pair
distortion. Nearby electrons feel force of attraction by the region of increased

charge density but this force of attraction is over powered by the force of mutual
repulsion between nearby electrons. However, electrons far away get bound
with each other as the force of mutual repulsion between distant electrons is
very small and force of attraction due to charge distortion over rides the repul-
sion. In this way Cooper electron pairs are formed. This is the reason why the
distance between the two electrons of Cooper pair may range from 50 to 100 nm
or more. This is in comparison to the lattice separation; distance between two
neighbouring ions of the lattice, of 0.1–0.4 nm. It is important to note that only
those electrons that are far apart may form Cooper pairs. Figure 2.45 shows
the quantum mechanical interaction between the two electrons of the Cooper
pair.
SAQ: What are phonons? What role do they play in the formation of Cooper
pairs.
SAQ: Cooper pairs are formed with far away electrons, explain.
2.6.2 BCS Theory of Superconductivity
John Barden, Leon Cooper and J. Robert Schrieffer, in 1957 developed a quantum
mechanical microscopic theory for superconductivity, which in short is termed as
BCS theory. The theory explains the resistance less flow of current by paired electrons
in some materials below critical temperature T C . The theory is based on the concept
of Cooper pairs which are formed in superconducting specimen below the critical
temperature. The salient features of the theory may be summarised as follows:
• Phonon-(free) electron interactions in some materials, below critical temperature,
give rise to the formation of Cooper pairs.
• The binding energy of a cooper pair is very small; of the order of few milli-electron
volts (≈ 10−3 eV). Therefore, to keep Cooper pairs intact, the superconducting
specimen must be kept below the critical temperature T C .
• Electrons 100 nm or more far apart from each other join to form Cooper pair. It
is because, the force of repulsion between distant electrons is small and may be
overcome by the force of attraction between the phonon and electrons. In case of
nearby electrons the force of repulsion is high and dominates over attraction by
phonon.
• BCS theory requires that the linear momentum of the Cooper pair must be zero;
therefore, the two electrons forming a Cooper pair must be moving in opposite
direction.
• The bound state of two electrons as Cooper pair is lower in energy than the energy
state of unbound free electron, the state corresponding to Cooper pair lies below
the Fermi level.
• Formation of Cooper pair is a transient phenomenon. Suppose at a given instant
a Cooper pair has two electrons marked-1 and 2. It is possible that at the next
instant electron-1 of the Cooper pair may change its partner and another electron
out of the large number of free electron may join with electron-1 to form the
Cooper pair, at the next instant electron marked-1 may be replaced by some other
electron in the Cooper pair and so on. In this way, almost all free electrons get
coupled with each other in forming Cooper pairs. This results in inter-linking or
coupling of all free electrons in the volume of the superconducting specimen and
they move coherently.
• Superconductivity is a quantum mechanical phenomenon; each free electron in
quantum mechanics is represented by a wavefunction that extends over the volume
of the specimen. Wavefunctions of all free electrons therefore overlap generating
a resultant wavefunction representing all free electrons together. The resultant
wavefunction gives rise to the coherent behaviour of all free electrons.
• A Cooper pair has two electrons each with spin 1/2 ℏ and, therefore the spin of
a Cooper pair may either be 0 or 1 ℏ. Particles with integer spin (0, 1 ℏ, 2ℏ…)
are called Bosons and obey Bose–Einstein quantum statistics. Unlike electrons
which follow Fermi- Dirac statistics and not more than two electrons with their
spins in opposite directions can stay in an energy state, large number of Bosons
may stay in a given energy state. As such all Cooper pairs in the superconducting
specimen stay in the lowest energy state, the ground state. Simultaneous stay of
all Cooper pairs in the ground state is referred as Condensation.
• In energy band diagram of a superconductor, the ground band of Cooper pairs is
separated by a forbidden energy gap Eg from the energy band for free electrons, as
shown in Fig. 2.46. The forbidden energy E g is related to the transition or critical
temperature T c by the relation.
E g = 3.53kTC , here k stands for Boltzmann constant.
Energy band diagram for a metal is also shown in Fig. 2.46 for comparison.
• In the superconductive state the current flow is constituted by the motion of
coherent Cooper pairs. In ordinary metal resistance essentially develops out of
the inelastic collisions between the lattice and free electrons, however, in case of
superconductors, where current is constituted by the coherent motion of Cooper
pairs, inelastic scattering between lattice and coherent Cooper pairs is not possible.
It is because for inelastic collisions the Cooper pairs must change into free elec-
trons for which energy equivalent to forbidden energy is required. Below critical
Fig. 2.46 Band diagram for a normal metal b superconductor
temperature this energy is not available. Hence current constituting motion of

coherent Cooper pairs is free of scattering, without ant resistance.
SAQ: Stable levitation is more easily achieved with type-II superconductors. Why?
SAQ: In a bench-top levitation demonstration the permanent magnet is first pushed
towards the superconductor, why?
SAQ: What is meant by the transient nature of Cooper pair formation?
SAQ: Why Cooper pair-lattice collisions do not take place in the current flow
through a superconductor?
Problems
P2.1 Photons of frequency larger than 1 × 1015 Hz get absorbed by a semiconductor.

What is the Forbidden energy gap of the material?
ANS: 4.1357 eV
P2.2 A semiconductor has a forbidden energy gap of 1.46 eV. Calculate the
maximum wavelength of photons that will be absorbed by the material.
ANS: 0.85 × 10−6 m
P2.3 Copper has 8.5 × 1028 free electrons per m3 , calculate the valency of copper
given that its molecular mass is 63.54 g/mol and density 8.96 g cm−3 .
ANS: 1
P2.4 A metal at very low temperature near absolute zero has Fermi energy of 7.0 eV.
Calculate the number density of free electrons in the metal assuming that the
effective mass of electron is equal to its actual mass 9.1 × 10–3 1 kg.
ANS: 8.28 × 1028 m−3
SA2.1 Define a semiconductor and differentiate it with insulator

SA2.2 What is meant by the breakdown of an insulator? Define the dielectric
strength of an insulator.
SA2.3 Explain why the resistivity of a semiconductor decreases with temperature
while the resistivity of a conductor increases with temperature.
SA2.4 What is the underlying principle of zone refining technique? Give an outline
of steps to obtain ultrapure Silicon from its ore.
SA2.5 What is doping? Describe ion implantation technique of doping.
SA2.6 Draw energy band diagram and covalent bond picture of a p-type Silicon
semiconductor at T > 0 K, indicating Fermi level, acceptor level, etc.
SA2.7 What is depletion layer and how does it form in a p–n junction? What are
the special properties of depletion region?
SA2.8 In depletion layer the total uncovered positive charge on n-side is equal to
the total uncovered negative charge on p-side. Suppose in a p–n junction p-
side is heavily doped as compared to the n-side, explain why the depletion
layer will extend more on the n-side.
SA2.9 What is a Cooper pair? How is it formed? Discuss the meaning of
condensation with reference to superconductivity.
SA2.10 State salient features of BCS theory for superconductors
SA2.11 A superconducting body acts as a perfect diamagnetic material, explain.
SA2.12 How can one explain the absence of resistivity in superconductors?
SA2.13 What is Hall Effect? What is its significance? Define Hall coefficient and
write its units.
SA2.14 Drive an expression for Hall voltage.

Note: Some of the following multiple choice questions may have more than one
correct alternative. All correct alternatives must be marked for complete answer.
MC2.1 Conductivity has units
(a) Ω-m (b) Ω m−1 (c) Ω−1 m−1 (d) Siemen per metre
ANS: (c), (d)
MC2.2 Dielectric strength is
(a) The maximum electric fields that an insulator may withstand before
electric breakdown
(b) Maximum voltage that may be applied to an insulator before electric
breakdown
(c) Maximum current that may pass through a conductor before electric
breakdown
(d) Maximum voltage that may be applied to a conductor before electric
break down
ANS: (a)
MC2.3 If ni , np and ne , respectively, denote the charge carrier concentration
in intrinsic material, concentrations of holes in p-type material and the
concentration of electrons in n-type material, then
(a) n i > n p (b) n i > n e (c) n p > n i (d) n p n e = n i2
ANS: (c), (d)
MC2.4 Resistance of a piece of a semiconductor decreases with the increase of
temperature because;
(a) The mean free path of electron–lattice collisions increases with the rise
in temperature
(b) Rate of lattice-electron collisions decreases with the rise of temperature
(c) With the increase of temperature, the rate of increase of carrier
concentration overtakes the rate of rise of lattice–electron collisions
(d) With the increase of temperature both the rate of increase in carrier
concentration and rate of increase of lattice-electron collisions are nearly
equal
ANS: (c)
MC2.5 When dopant concentration in a given piece of a semiconductor is
increased, the resistance of the piece.
(a) Increases
(b) Decreases
(c) Remains unaltered
(d) May increase or decrease depending the material
ANS: (b)
MC2.6 At room temperature, the thermal velocity of free electrons in a semicon-
ductor is of the order of
(a) 10−5 m
s
(b) 100 m
s
(c) 105 m
s
(d) 1015 m
s
Ans: (c)
MC2.7 Zone refining technique is based on the principle that;
(a) Mobility of impurity atoms is less in molten state, and the melting
point of pure Silicon is lower than the impure Silicon
(b) Mobility of impurity atoms is more in molten state, and the melting
point of pure Silicon is lower than the impure Silicon
(c) Mobility of impurity atoms is less in molten state, and the melting
point of pure Silicon is higher than the impure Silicon
(d) Mobility of impurity atoms is more in molten state, and the melting
point of pure Silicon is higher than the impure Silicon
ANS: (d)
MC2.8 Depth profile of implanted impurity diopant ions is
(a) Uniform (b) Shows a single peak at the range of implanted ion (c)
Uniformly decreases (d) Uniformly increases
ANS: (b)
MC2.9 Given that d p and d e , respectively, represent the width of depletion layer
on the p-side, and on the n-side while N p and N e, respectively, the
concentration of dopants on the n- and p-sides and if d p = 1.5 d n , then;
(a) N p = N e (b) N p = 3 N e (c) N p = 1.5 N e (d) N e = 1.5 Np
ANS: (d)
MC2.10 Reverse current in a pn junction gets doubled for every 100 C rise of
temperature because;
(a) Reverse current is constituted by minority carriers that are produced
by bond breaking which increases with temperature
(b) Reverse current is constituted by majority carriers that are produced
by bond breaking which increases with temperature
(c) Reverse current is constituted by minority carriers that are produced
by bond breaking which decreases with temperature
(d) Reverse current is constituted by majority carriers that are produced
by bond breaking which decreases with temperature
ANS: (a)
MC2.11 Which of the following always has zero magnitude in a superconducting
material?
(a) Electric field E (b) Potential difference V (c) Current I (d) Magnetic
susceptibility χ
ANS: (a), (b)
MC2.12 In an experiment a ring shaped specimen of superconducting material
is placed in a magnetic field of strength B at temperature T > T C . The
temperature T is then reduced and maintained below T C . The magnetic
field is switched off. Which of the following statement(s) will correctly
describe the result of the experiment?
(a) Magnetic field will remain trapped outside the ring opening but it will
be zero inside the ring opening
(b) Magnetic field will remain trapped inside the ring opening and will be
zero outside the ring opening
(c) Magnetic field will be zero both inside and the outside the ring opening
(d) Magnetic field will remain trapped both inside and outside the ring
opening
ANS: (b)
MC2.13 A current I x in positive X-direction passes through a slab of n-type semi-
conductor. If a magnetic field BZ in positive Z-direction is applied across
the semiconductor slab, the deflection force due to magnetic field on free
electrons constituting the current will be in the direction of
(a) X-axis (b) Y-axis (c) Z-axis (d) 45° between X- and Z-axis
ANS: (b)
MC2.14 Units for Hall coefficient are
m3 C m3 A.s
(a) C
(b) m3
(c) A.s
(d) m3
ANS: (a), (c)
LA2.1 Discuss the electron energy band theory of crystals and hence explain the
classification of solids according to their electrical properties.
LA2.2 What are the characteristics of Insulators? Why do they have very large
value for resistivity? Draw sketches for the energy band pictures of an insu-
lator and a conductor. Explain the phenomena of breakdown in insulators
and define the dielectric strength.
LA2.3 Draw labelled diagrams for the band structures of a p-type and an n-type
semiconductor. What are acceptor and donor levels and how do these levels
affect the Fermi level?
LA2.4 With the help of a labelled energy band diagram discuss the formation of a
pn junction. What is meant by thermal equilibrium and how it is achieved in
case of a pn junction? Discuss the formation of depletion layer at junction
boundary and list some of its properties.
LA2.5 What is a pn junction diode? Discuss the current flow through a junc-
tion diode under forward and reverse bias and draw its current voltage
characteristics.
LA2.6 With necessary details describe ion implantation method of doping of
semiconductor wafer.
LA2.7 What is Meissner effect? How can it be explained? Discuss stable levitation
and essential conditions to achieve it.
LA2.8 What are Cooper pairs and their condensation? Outline BCS theory of
superconductivity and explain why lattice-Cooper pair collisions do not
take place in superconductors.
LA2.9 Give distinguishing features of type-I and type-II superconductors. Wires
of type-II superconductors are often used for making strong electromag-
nets, explain.
LA2.10 What is isotope effect and why it was so significant in developing a
theory for superconductivity? Explain the process of Cooper pair formation
via electron–phonon interaction. With the help of energy band diagrams
explain the difference between a conductor and a superconductor.
LA2.11 What is Hall Effect? Describe with the help of a diagram the setup for
measuring Hall voltage and derive an expression for Hall voltage in terms
of Hall coefficient.
Chapter 3
Magnetic Materials
Objective
Origin of magnetism in matter, types of magnetisms, their properties, applications,
etc. will be discussed in this chapter. It is expected that after reading this chapter the
reader will be able to understand how magnetic properties develop in materials and
how materials may be classified in terms of their magnetic behaviour. He will also
learn how materials with desired magnetic properties may be developed.
3.1 Introduction
Magnetic materials include a variety of materials that are used in diverse applica-
tions. It is interesting to note that magnetic materials are utilised in generation and
distribution of electricity and in most cases, they are also used in the appliances that
use that electricity. Magnetic materials are used for the storage of data on audio and
video tapes and computer disks. Magnetic materials also have applications in the
field of medicine, they are used in body scanners as well as a range of applications
where they are attached to or implanted into the human body. Non-polluting electrical
vehicles have very efficient motors that utilise advanced magnetic materials.
The fact that Earth behaves like a bar magnet was known to Indian saints and seers
many centuries ago, and they developed and devised codes called ‘Vastushastra’ for
build temples and buildings based on Earth’s magnetic meridian.
In the modern era, however, in 1600 William Gilbert published the first systematic
experiments on magnetism in the pamphlet ‘De Magnet’. By the end of the eighteenth
century, scientists have noticed many electrical and magnetic phenomenon but they
all believed that these two branches of science, the electricity and the magnetism
are quite independent/separate from each other and that perhaps there is no direct
relationship between the two. Lightning, the electric phenomena, was well known to
ancient people but first magnetic material that showed the power of attracting small
iron pieces was the loadstone, mineral magnetite (Fe3 O4 ), found in form of small
https://doi.org/10.1007/978-3-031-32084-2_3
138 3 Magnetic Materials
pieces of rock, perhaps in Greece. Then in July 1820, Danish natural philosopher
Hans Christian Orsted published a pamphlet that showed clearly that electric current
and magnetism are very closely related to each other. It is said that Orsted, born in
August 1777 after completing his Ph.D. in philosophy in 1801, travelled through
Europe, as was customary at that time, and met many scientists/philosophers of
Germany, France, etc. One person he met was Johann Ritter, a scientist who believed
that electricity and magnetism are related. Orsted might have been influenced by
him.
Orsted returned back to Copenhagen in 1803 and tried for a faculty position at
University there, but initially did not succeed. However, he started private lecture
courses, charging admission fee, which attracted many young people. Later he got
a position in Copenhagen University, perhaps on account of his popularity through
private lectures.
Though most of the scientists at that time thought electricity and magnetism to
be two totally independent attributes of matter, but there were some indications, for
example it was observed that a magnetic compass if struck by lightning reverses its
poles; that pointed to some connection between the two. From his writings/lectures
it appears that Orsted believed that electricity and magnetic behaviour are two inde-
pendent properties of all matter and that these properties might interfere with each
other.
During one of his lecture demonstration on 21 April 1820 while setting up his
apparatus, Orsted noticed that whenever he switched on current in his circuit, the
north pole of the magnetic needle placed nearby deflected a bit. He also noticed that
the direction of deflection of the North Pole of the compass needle changed when the
direction of current in his circuit was reversed and that the effect of current on the
movement of magnetic needle may be shielded by interposing an insulator/dielectric
material between the electric circuit and the magnetic needle. Orested published his
findings which were mainly qualitative, in a pamphlet, circulated to other scientists,
on 21 July 1820; but the effect was clear: An electric current generates a magnetic
field.
French scientist Andre-Marie Ampere opined that the production of magnetic
field by electric current is the fundamental feature of magnetism. In fact, he went to
the extent that all magnetism, permanent magnets including, is the result of currents
in them. Later, some 160 years ago, in 1864, James Clerk Maxwell carried out the
first profound unification of Nature’s two forces, the electric force and the magnetic
force in form of his famous Maxwell’s equations.
The present understanding is that stationary electric charges produce an elec-
tric field around them; charges moving with uniform speed (current elements) give
rise to both the electric and the magnetic fields while accelerated charges radiate
electromagnetic fields.
3.2 Electric Current and Magnetic Field 139
3.2 Electric Current and Magnetic Field
Electric current produces a magnetic field which may be visualised as a pattern of

circular magnetic field lines surrounding the current carrying wire. The direction of
the magnetic field may be determined using the compass needle while a Hall probe
may be used to determine the magnitude of the field. Careful experimental study
of the direction and magnitude of the magnetic field produced by a current in an
infinite (very long) straight wire reviled the right-hand rule, according to which if
one aligns the thumb of the right hand with the direction of current flow in the straight
wire, the curled fingers of the right-hand point in the direction of the magnetic field.
The magnitude of the magnetic field strength B due to current I in an infinite (or a
sufficiently long) straight conductor at a perpendicular distance r from the conductor
is given as;
μ0 I
B= (3.1)
2πr
Here constant μ0 is the permeability of the free space, is a basic constant of nature
related to the velocity of light c, having the value μ0 = 4π × 10−7 T m/A. Since the
current carrying conductor is very long, the magnetic field strength B at the point
of observation O depends only on the perpendicular (shortest) distance r from the
conductor.
Biot–Savart argued that each little segment dl of the current produces a magnetic
field at the point of observation and that the total magnetic field due to the complete
current carrying conductor (of any shape) is given by,
∫ −
→ ∫ −
→
μ I dl X r̂ μ0 I dl × →r
B→ = 0 = (3.2)
4π r 2 4π |r|3
−
→
In Eq. (3.2) dl is the vector element of length in the direction of the current flow
and r̂ is the unit vector in the direction of vector distance r from dl to the point
of observation O. The line integration is to be carried on the length of the current
carrying conductor.
Figure 3.1a shows the direction of magnetic field lines due to an infinite straight
conductor carrying current I, Fig. 3.1b depicts the right-hand rule to specify the
direction of magnetic field lines, Fig. 3.1c gives the magnitude of magnetic field
strength at point O situated at a perpendicular distance r from the straight infinite
conductor carrying current I and Fig. 3.1d shows the strength of magnetic field due
to a current element dl at distance r from it as given by Biot–Savart law.
A bar magnet suspended in Earth’s magnetic field orients itself in North–South
direction. The North seeking end of the bar magnet is called the North end and the
geographic South seeking end as the South end. When lines of magnetic field of a bar
magnet are drawn using a compass needle, they appear to originate from a point on
the North end and appear to terminate at a point on the South end. Since all magnetic
Fig. 3.1 a Lines of magnetic field of an infinite conductor carrying current. b Right-hand rule giving
the direction of magnetic field. c Magnitude of magnetic field strength at a point O at distance r from
an infinite conductor carrying current I. d strength of magnetic field due to current element dl at
point O
field lines appear to originate from one point on the North end, this particular point
is called the North Pole of the magnet, and similarly point on the south end where
all magnetic field lines appear to terminate, the South Pole. Number of magnetic
field lines per unit area around a point is a measure of the intensity of magnetic
field at that point. It is obvious that the magnitude of magnetic field intensity at
poles is a maximum. Figure 3.2a shows the magnetic field lines of a bar magnet
drawn using a compass needle. Though when looked from outside it appears as if
the lines of magnetic field of a bar magnet originate from North pole and terminate
at South pole, but the fact is that magnetic field lines are continuous within the bar
magnet as shown in Fig. 3.2b. Parallel magnetic field lines at some place indicate a
uniform magnetic field in that region. North poles as well as the south poles of two
bar magnets repel each other while the opposite poles attract. One special property
of a magnet is that North and South poles occur in pairs, for example if one break a
bar magnet into pieces, each piece develops North and South poles. It is to say that
free North or South magnetic poles (mono poles) never occur in nature. Magnetic
poles to some extant may be compared with electric charges, magnetic north pole
like positive electric charge and magnetic south pole like negative electric charge.
3.3 Magnetic Dipole Moment 141
Fig. 3.2 Magnetic field lines of a bar magnet
However, there is one fundamental difference; it is possible to have an isolated single

free positive or negative electric charge but it is not possible to have a free single
magnetic north pole or a single free magnetic south pole. Magnetic North and South
poles always occur in pair. A major point of difference between the electric and
magnetic field lines is that electric field lines actually originate from positive charge
and terminate at negative charge but magnetic field lines always form a closed loop.
Though electric and magnetic fields are two quite different fields, however, it
is observed that the electric field due to an electric dipole, (two equal and opposite
charges separated by a small distance) is similar to the magnetic field of a bar magnet
as may be seen in Figs. 3.2 and 3.3a.
Lines of magnetic fields due to circular current loop are shown in Fig. 3.3b, c. It
may be observed in these figures that the magnetic field due to a bar magnet (or a
magnetic dipole) is similar to the magnetic field produced by a circular current loop.
Therefore, a circular current loop is also called a magnetic dipole.
SAQ: What information does the permeability of a material give about the magnetic
behaviour of the material?
3.3 Magnetic Dipole Moment
Magnetic dipole moment or simply magnetic moment of an object is a vector quantity

used to measure the tendency of the object to interact with an external magnetic field
and is represented by the symbol µ. The object’s intrinsic magnetic properties play
an important role in deciding the tendency of its interaction with external magnetic
field. The intrinsic magnetic properties of the object are often visualised as emanating
from a tiny bar magnet with north and south poles, hence the name dipole moment.
Figure 3.4a shows the magnetic dipole moment for a small bar magnet of length
2a and pole strength M when looked from a distance. A current carrying circular
loop of area A (= πr 2 ) when looked from a distance produces a magnetic field that
resembles the magnetic field of a tiny bar magnet and the magnetic (dipole) moment
of the current loop is equal to IA, where I is the current in the loop and A is the area
Fig. 3.3 a Electric field lines (blue) and equipotential lines (red) due to an electric dipole b magnetic
field lines due to a current loop with axis along Z-axis; the yellow line shows the segment of the
current loop c current loop and associated magnetic field lines
enclosed by the current loop. The direction of the magnetic moment of the loop may
be easily determined using the right-hand rule.
When an object having a magnetic dipole moment μ → is placed in an external
−−→
magnetic field B ex t , the object or the magnetic moment associated with the object
experiences a torque that tries to align the dipole moment of the object in the direction
of the external magnetic field. The magnitude and direction of the torque τ are given
−−→ −−→
by the vector equation τ→ = μ×→ B ex t as shown in Fig. 3.4c. In case when B ex t has the
magnitude of 1, and the angle between magnetic moment μ and external magnetic
field Bext is 90°, then τ = μ. One may, therefore, define the magnetic moment of an
object as the maximum torque experienced by the object in unit magnetic field.
3.4 Magnetic Moment of a Charged Particle Moving in a Circular Orbit 143
Fig. 3.4 a Magnetic dipole moment of a small bar magnet. b Magnetic dipole moment of a circular
current loop. c Torque experienced by a magnetic dipole moment in an external magnetic field.
d Potential energy of a magnetic dipole in an external magnetic field
The potential energy U of the interaction between the magnetic moment and the
−−→
external magnetic field is given as U = −μ.→ B ex t . The negative sign in expression for
U is included to indicate that the potential energy of interaction will be a minimum
when both the magnetic moment and the external magnetic fields are parallel and
will be a maximum when they point in opposite directions. It is arbitrarily chosen
to assign a value zero to the potential energy U when magnetic moment points in a
direction perpendicular to the external magnetic field.
3.4 Magnetic Moment of a Charged Particle Moving

in a Circular Orbit
Figure 3.5 shows a particle of mass m, charge + q moving with linear velocity v
in a circular path of radius r. We will calculate the magnetic dipole moment µ and
angular momentum L of the particle under classical approach.
Let us first calculate the magnetic (dipole) moment μ → of the particle. A particle
having an electric charge + q, moving in a circular closed path with linear velocity
v (angular velocity ω = v/r ) crosses a fixed point on its circular path after each
Fig. 3.5 Magnetic dipole

moment and angular
momentum of a particle of
mass m and charge q moving
with linear velocity v in a
circular path of radius r
time interval t = 2πr/v. The frequency of crossing the fixed point f = 1t = v/2πr
and the current I constituted by the circulating charge + q and defined as the rate of
flowing of charge, may be written as,
+qv
I = +q f = (3.3)
2πr
Current I flows in the same direction as the direction of motion of the electric
charge + q, in case the charge q is negative, the current I will flow in the direction
opposite to the direction of motion of the charge. The circular area A enclosed by
current I is A = π r 2 and, therefore, the magnetic (dipole) moment of the current loop
is given by,
−
→ +qv ( 2 ) +qvr
μ| = IA = πr = (3.4)
2πr 2
Next, let us calculate the angular momentum L of the charged particle of mass m
which is moving in the circular path of radius r with linear velocity v. By definition,
the angular momentum is given as;
| |
| |
L→ = r→ × p→ = r→ × (m . v→) or | L→ | = mr v sin θ (3.5)
In Eq. (3.5) θ is the angle between the linear velocity v and the radius vector r,
which in the present case of circular motion is always 90° and hence
→ = mrv
| L| (3.6)
The direction of L is given by the vector cross-product rule and in the present
case angular momentum of the particle points in the upward direction, normal to the
−
→
plane containing r and v. It may be observed that if q is positive, both −
→
μ and L
points are in the same direction but if q is negative, −
→
μ will point downwards while
angular momentum L will point upwards.
3.4 Magnetic Moment of a Charged Particle Moving in a Circular Orbit 145
Having obtained expressions for the magnetic moment of the moving charge q
and its angular momentum we calculate the ratio,
| | +qvr
|μ |
|→|=γ = 2 = q (3.7)
| L→ | mr v 2m
The above ratio of magnetic moment to angular momentum of a particle is called

the Gyromagnetic ratio, represented by the symbol γ , and is an important parameter
for microscopic particles. Equation (3.7) tells that units for γ are Coulomb per
kilogram in SI system. Also, one may write,
q →
→ = γ L→ =
μ L (3.8)
2m
3.4.1 Classical to Quantum Mechanics
So far we considered the motion of a charged particle in classical limits and derived
expressions for magnetic dipole moment and angular momentum. However, it is
known that the behaviour of microscopic particles like neutron, proton and elec-
tron, etc., which are constituents of atom, is better understood in terms of quantum
mechanical treatment. It is interesting to note that all classical expressions derived
here holds good in quantum mechanical treatment also. However, there are two points
of difference between the classical and the quantum mechanical treatments: (i) In
quantum treatment it is not important wether the current flows in a circular path or
not, current must flow in a closed loop enclosing some area A (which may not neces-
sarily be a circular area) and that value of area may be used in above expressions. (ii)
In case of classical treatment L may have different continuous values from zero to
mv, depending on the value of angle θ, but in quantum physics L may assume only
discrete values. In switching from classical to quantum treatment, the expression for
L becomes,
√
L = l(l + 1)ℏ (3.9)
where l is either 0 or a positive integer, i.e. l may have values 0, 1, 2, 3, … and

ℏ is the unit in which angular momentum is measured in quantum mechanics ℏ =
1.05457 × 10−34 J s.
‘l’ is called the orbital angular momentum quantum number and the state of
motion of the particle is characterised according to the value of l; when l = 0
the state of relative motion is called the s-state, for l = 1, the p-state; for l =
2, the d-state; for l = 3, the f-state and so on.
Thus in quantum mechanical treatment the magnetic dipole moment of a charged
particle for orbital motion may be written as;
−
→l q →
μ = L (3.10)
2m
and,
√
L= l(l + 1)ℏ (3.11)
where l may have values 0, 1, 2, 3, …; and ℏ = 2π h

, h being Planck’s constant.
Further in quantum mechanics it is convenient to measure angular momentum in
units of ℏ.
3.5 Magnetic (Dipole) Moment of Electron
Electrons are constituents of all atoms and molecules. It is known that two kinds of
motions, namely orbital motion around the atomic nucleus and inherent spin motion,
are associated with electrons. Both these motions obey laws of quantum mechanics
and have corresponding magnetic moments associated with them. Let us first consider
the orbital motion of an electron.
(i) Orbital motion of electron
The magnetic moment μorb e of electron due to its orbital motion may be obtained
from Eq. (3.10) by replacing (i) q by (− e) where minus sign indicates that the
charge of electron is negative, e = 1.6 × 10−19 C is the charge of the electron
and (ii) m by m e , the mass of electron. With these substitutions, we get
−−→ e →
μorb
e =− L (3.12)
2m e
Equation (3.12) tells that in case of electron, the direction of magnetic moment
μorb
e is opposite to the direction of angular momentum L, because of the negative
charge on electron.
The simplest way of measuring the magnetic (dipole) moment of a particle is
to put the particle in some external field in a given direction and measure the
interaction energy U, from which the magnitude of magnetic moment may be
calculated. Suppose the external magnetic field is applied along the z-axis. Then
the z-component of the magnetic moment will come into play. It is reasonable
to assume that Eq. (3.12) will also hold for the z-components of μle and L,
therefore,
( ) e
μorb
e z
=− Lz (3.13)
2m e
3.5 Magnetic (Dipole) Moment of Electron 147
Now quantum mechanics tells that the z-component of angular momentum is

also quantised, i.e.
L z = ml ℏ (3.14)
Here, ml , called the magnetic quantum number,

may have (2l + 1) different values, starting from
−l, (−l + 1), (−l + 2), (−l + 3), . . . , (l − 3), (l − 2), (l − 1), l. For
example if orbital angular momentum l = 0 (s-state), then magnetic quantum
number m l may have only one value that is 0. If l = 1(p-state), magnetic
quantum number m l may have three different values; − 1, 0, and + 1; if
l = 2(d-state) then m l may have 5 different values, − 2, − 1, 0, + 1 and + 2
and so on.
Substituting the value of L z from Eq. (3.13) in Eq. (3.14), one gets;
[ ]
( ) e e eℏ
μorb
e z
=− Lz = − ml ℏ = − ml (3.15)
2m e 2m e 2m e
The negative sign on the RHS of Eq. (3.15) simply indicates that the direction
of μorb
e is opposite to L z .
[ ]
eℏ
Quantity 2m e
is defined as 1 Bohr Magneton represented either by Bm or
μB . Magnetic moment of atomic particle like electron is measured in units of
Bohr magneton. Magnitude of 1 Bm or 1μB depends on the system of units. In
SI system of units 1 Bm = 2meℏ
e
while in Gaussian CGS system 1 Bm = 2meℏe c
Metric (SI) Equivalent 9.274 × 10−24 J/T or m2 A

CGS 9.274 × 10−21 erg/G
eV 5.7883 × 10−5 eV/T
Bohr magneton may be expressed by different dimension formulas depending

on systems of units, some of them are;
M L 3 T −1 Q −1 ; N m( A/m)−1 ; N M T −1 etc.
As has been mentioned, in experimental measurement of the magnetic moment

of a particle, an external magnetic field is applied in a certain direction (denoted
by direction z) and the z-component of the magnetic moment is measured. Thus,
any measurement of magnetic moment yields only the component of magnetic
moment in a specified direction. Further, the measurement may not give a
unique value of the component; it is because, from Eq. (3.15) the measurement
may give one of the several possible values of magnetic quantum number
m l . Let us understand it by taking an example. Suppose one measures the

component of magnetic moment of an electron due to its orbital motion in state
− d where l = 2. The possible values of m l for the case (l = 2) are: − 2, −
1, 0, + 1 and + 2. The measurement may give any one of these values. That
means when same measurement is carried out by different experimentalists or
repeated, they may get different values for the same component of magnetic
moment. To remove this ambiguity, it has become a convention to (i) treat
the measured component of magnetic moment as the magnetic moment of the
particle (not to call it a component) and (ii) out of the several possible values
of m l , the maximum positive value is taken as the measured value. Suppose in
a given case J is the maximum positive value of magnetic quantum number,
then one writes;
[ ]
eℏ
μorb
J = − J (3.16)
2m e
As already mentioned the negative sign in Eq. (3.16) may be dropped, taking in
account the fact that the magnetic moment of electron due to its orbital motion
is opposite in direction( to the
) angular momentum. It follows from Eqs. (3.15)
and (3.16) that when μorb e z
is measured in units of Bohr magneton (Bm ) and
proper sign of μorb
J is taken into account, then the ratio
( orb )
μe z (in unit Bm ) μorb
= 1 or J
=1 (3.17)
ml J
Expression (3.17) may be generalised for any quantised magnetic dipole

momentum (measured in Bohr magneton) and the corresponding angular
momentum quantum number in the following form;
μ(in Bm )
=g (3.18)
J
Equation (3.18) implies that if μ is the value of the magnetic moment of some
particle measured in units of Bohr magneton (Bm ) and J is the corresponding
quantum number associated with motion of the particle, then the ratio of the
two may be represented by a constant g. The constant g is called the g-factor of
the particle and depends on the nature of the particle as well as on the type and
state of motion of the particle. For example, the value of g-factor for electron
for its orbital motion georbit may be given as;
μorb
J
= georbit = 1 (3.19)
J
3.5 Magnetic (Dipole) Moment of Electron 149
Thus, the g-factor (a dimension less constant) of electron for its orbital motion
georbit = 1.
Difference between Gyromagnetic ratio γ and the g-factor: Both represent
the ratio of magnetic moment to angular momentum, but units for measuring
the magnetic moment and the angular momentum are different. The ratio of
magnetic moment to angular momentum quantum number when magnetic
moment is measured in units of Bohr magneton, gives the g-factor.
(ii) Spin motion of electron
Motion of electron around the nucleus of the atom is termed as its orbital motion.
Apart from that an electron also possesses an inherent motion called the spin
motion. Spin or spin motion is purely a quantum mechanical concept that cannot
be explained in terms of classical physics. It is reasonable to assume that the
inherent spin motion of electron follows the same laws as are followed by its
orbital motion. There is an angular momentum quantum number L associated
with the orbital motion similarly there is a quantum number S associated with
spin motion. One can write down an expression for spin motion corresponding
to Eq. (3.18) of orbital motion as given below;
spin
μs (in units of Bm )
= g spin (3.20)
s
Or
μspin
s (in units of Bohr magneton (Bm )) = gespin s (3.21)
spin
Precise measurement of spin magnetic moment μs of electron gives the value
spin
of ge for electron as 2.002318. Therefore, in case of electron;
georbit = 1 and gespin = 2.002318
spin
The fact that ge ≈ 2, while georbit = 1 implies that spin motion of electron is
twice as effective in producing magnetic field as the orbital motion.
There are many electrons in an average atom and the total magnetic dipole
moment due to all electrons in the atom may be calculated in two different
ways: (a) In this method, called the j − j coupling method, the first step is to
find the total angular momentum ji of the ith electron by quantum mechanically
adding its orbital and spin angular momentums, ji = li + si . In the second step
the total angular momentum J of all the electrons is determined by quantum√
mechanical adding of total angular momentums ji ’s, of electrons; J = ji .
Once total angular momentum J is known one may calculate the magnetic
moment of the atom due to all its electrons using the expression
( j− j) ( j− j )
μatom = gatom J (3.22)
( j− j)
Here gatom is the g-factor for j − j coupling. It may be mentioned that quantum
mechanical addition of angular momentums is quite different from the classical
vector addition. (b) In second method, called (l − s) coupling, the total orbital
angular momentum L of all electrons in the atom is determined √ by quantum
mechanical adding of their orbital angular momentums, L = li and simi-
larly total spin S of the atom is√
obtained by quantum mechanical adding of spins
of individual electrons, S = si . The resultant angular momentum J of the
atom is then obtained by quantum mechanical addition of L and S; J = L + S.
Expression μ(L+S) (L+S)
atom = gatom J may be used to obtain the magnetic moment of
(L+S) ( j− j )
the atom. g-factors; gatom and gatom may have different values for the same
atom.
(iii) Magnetic moments of nuclear particles
Nucleus of every atom contains protons and neutrons which also perform orbital
and spin motions. Nucleons (both neutron and proton), therefore, like electrons,
have magnetic moments. Magnetic moment of nucleons is measured in a unit
called nuclear magneton, denoted by Nm . 1 Nuclear Magneton (Nm ) = 2m eℏ
p
,
where m p is the mass of proton. One nuclear magneton is about 1836 times
smaller than one Bohr magneton because a proto (1.6 × 10–27 kg) is roughly
1836 times heavier than an electron (9.1 × 10–31 kg). As a result, the contri-
bution of magnetic moments of nucleons is neglected while discussing the
magnetic moment of the atom.
As a concluding remark it may be said that magnetic fields originate from
electric currents.
3.6 Magnetic Behaviour of Solids
A common man may categorise solid matter in two classes: either magnetic or non-
magnetic. Magnetic materials are those which attract some other materials like iron
filing, lodestone, pins, etc., whereas non-magnetic solids, like wood, common salt,
chalk, etc. do not show any effect of magnetic field. However, the fact is that the
so-called non-magnetic materials are also affected by magnetic fields but the effect
is so weak that very sensitive detecting instruments and high magnetic fields are
required to demonstrate the effect.
Magnetic behaviour of matter originates essentially from the spin and orbital
motions of electrons in atoms/molecules or ions. Fortunately, all electrons in the
atom or ions, etc. do not contribute to magnetic properties. According to quantum
mechanical laws, an even number of electrons in a given energy level of the atom/ion,
etc. orient their orbital and spin motions in such a way that their magnetic moments
cancel out, resulting in no net magnetic field. Therefore, atoms, ions and molecules
the outer valence shell of which have even number of electrons do not show magnetic
3.6 Magnetic Behaviour of Solids 151
properties. However, if there are odd number of electrons in the outer energy levels of
an atom, ion or molecule, it has a magnetic dipole moment (and hence show magnetic
properties) that may be attributed to the unpaired electrons.
Large number of atoms, like that of bismuth, mercury, silver, copper, inert gases,
etc. have closed shell electron structure with even number of electrons, and hence,
do not show any magnetic properties of their own. However, when put in an external
magnetic field, these atoms develop an induced magnetism opposite to the applied
magnetic field. Such materials are called diamagnetic materials. Copper atom in
ground state has electron configuration: 1s2 2s2 2p6 3s2 3p6 4s1 3d10 ; and is diamag-
netic. Copper ion Cu1+ on loosing 4s1 electron also becomes a diamagnetic ion. Some
other examples are;
Copper (Cu) ground state: [Ar] 4s1 3d10 ; Copper ion (Cu1+ ):[Ar] 3d10 ; Lead (Pb)
ground state: [Xe]6s2 4f14 5d10 6p2 ; Silver (Ag) ground state: [Kr] 4p6 5s1 4d10 ;
Mercury (Hg) ground state: [Xe] 4f14 5d10 6s2
Atoms/ions with odd number of electrons in outer shell show magnetic properties
that may be attributed to the orbital and spin motions of the unpaired electron; such
atoms behave like small dipoles. Since atomic dipoles, in general, are randomly
oriented, the bulk material may show inherent magnetism but quite weak in strength.
These materials are classified as paramagnetic materials. In some special cases, tiny
atomic (or ionic or molecular) dipoles may align themselves in different patterns
giving rise to ferromagnetism, ferrimagnetism and antiferromagnetism.
In order to study different types of magnetism in details we define the
following important magnetic parameters that are frequently used to specify types
of magnetism.
SAQ: It is customary to neglect magnetic effects produced by the motion of protons
present in the nucleus of the atom, while calculating magnetic dipole moment
of the atom. Give justification.
SAQ: Neutron possesses a magnetic dipole moment due to its spin motion which
has negative value. What inference do you draw about the charge distribution
within a neutron?
3.6.1 Magnetic Induction B and Magnetic Field H
Magnetic induction may be defined in different ways; however, we describe here an

experimental method of determining magnetic induction B at a point in a magnetic
field. Suppose there is a magnetic field in free medium (vacuum or air) produced,
say, by a permanent magnet (see Fig. 3.6). It is known that conductor of length dl
carrying a current i and placed in this magnetic field at some point P will experience
a force F, given by;
F = B i dl sin θ , (3.23)
Fig. 3.6 Defining magnetic

induction at a point P in a
magnetic field
Here θ is the angle between the direction of current and the force.
In case of unit current (i = 1), force perpendicular to the direction of current per
unit length (dl = 1) is equal to the magnetic induction B at point P. Direction of
B will be perpendicular to the direction of current flow. The SI units for magnetic
induction B is tesla, represented by letter T. One tesla 1 T = one weber per square
metre that corresponds to 104 gauss. 1 T may also be represented as one kilogram
per second squared per ampere (kg s−2 A−1 ).
Having defined B at point P in terms of the force F, one may now define the
strength of magnetic field H at point P as,
B
H= (3.24)
μ0
Here μ0 is the permeability of the medium that is the( Hvacuum

) or air in the present case.
The value of μ0 for air or vacuum is 4π 10−7 Henery
meter m
or Newton per Ampere square
(N/A2 ). It may be noted that symbol μ is frequently used in science to denote different
quantities, for example in optics μ is often used to indicate the refractive index, it also
denotes the magnetic dipole moment, but here it is being used to represent another
magnetic parameter of a medium called permeability. Permeability of a medium is
a measure of the ease with which a magnetic field may be established in the medium.
The SI units of magnetic field strength H is amperes per metre (A/m), while the
CGS units are Oersted (Oe).
Let us now consider the case when the medium surrounding the point of obser-
vation P is not vacuum or air but is some other material, say, an iron or brass sheet
or any other material. The magnetic field H which is due to the permanent magnet
at point P will induce some magnetic field at point P because of the magnetisation
of the medium. If M denotes the induced magnetic dipole moment per unit volume
3.6 Magnetic Behaviour of Solids 153
Table 3.1 Permeability and

Material/medium Permeability μ (H/ Relative
relative permeability of some
m) permeability μr
materials
Vacuum/air 1.256637 × 10–6 1.000000
Copper 1.256629 × 10–6 0.999994
Iron (99.8% pure) 6.3 × 10–3 5000
Nickel 1.25 × 10–4 – 7.54 100–600
× 10–4
Superconductors 0 0
of the surrounding medium, then the magnitude of magnetic induction at point P (in
presence of the new medium) is given as’
B = μ0 (H + M) = μH (3.25)
Here, μ is the permeability of the medium (iron or brass, etc.). The ratio μμ0 =
μr is called the relative permeability of the medium is just a number/fraction
without units. Permeability and relative permeability of some materials are given in
Table 3.1.
In an isotropic medium B and H are parallel and the permeability is a scalar
quantity. However in an anisotropic medium B and H may not necessarily be parallel
and the permeability is a tensor.
The induced dipole moment per unit volume M, also called the intensity of
magnetisation or simply magnetisation is the effect of the magnetising field H on the
medium. Thus H is the ‘cause’ and M is the ‘result’. It is obvious that the ‘result’ M
will be proportional to its ‘cause’ H, i.e.
M ∝ H or M = χ H (3.26)
The constant of proportionality denoted by Greek letter χ is called the magnetic

susceptibility of the medium. Susceptibility χ is a parameter that demonstrates
the type of the magnetic material and the strength of the induced magnetic field in
the medium. Sometimes one uses mass susceptibility, denoted by χm , that may be
obtained by dividing χ with the density ρ of the medium. Further, it follows from
Eqs. (3.25) and (3.26) that;
B = μH = μ0 (H + M) = μ0 (H + χ H) = μ0 H(1 + χ )
Or
μH = μ0 H(1 + χ )
That gives;
μ r = (1 + χ ) (3.27)
Magnetic properties of different materials may be classified in terms of the sign

and the magnitude of its susceptibility χ . Since units of M and H are same, bulk
susceptibility χ is just a number/fraction and has no units and dimensions. However,
the mass susceptibility χm has units of metre3 /kilogram in SI system.
3.7 Classification of Magnetic Materials
Magnetic behaviour of all materials, depending on the sign and magnitude of their
magnetic susceptibilities, may be classified into five categories: (a) diamagnetic, (b)
paramagnetic, (c) ferromagnetic, (d) antiferromagnetic and (e) ferrimagnetic. Out of
some 90 stable elements of the periodic table around 38 are diamagnetic, 48 para-
magnetic, 03 ferromagnetic and only 01 antiferromagnetic at room temperature. No
element in natural form is found to exhibit ferrimagnetisms; only some compounds
such as mixed oxides show this type of magnetic behaviour.
3.7.1 Diamagnetic Materials
Materials with bulk susceptibility χ having small negative value in the range of
−10−6 to − 10−5 are classified as diamagnetic. Negative value of susceptibility
means that in an applied magnetic field diamagnetic materials acquire a magnetisation
which points opposite to the magnetising field H. Atoms of diamagnetic materials
do not show any magnetism in absence of any external magnetic field (H); that
essentially means that the outer shells of these atoms have even number of electrons
that chancel out their magnetic moments (due to their spin and orbital motions) in
pairs. However, when some external magnetic field H is applied to a diamagnetic
atom, it alters the speeds of rotation of the two electrons of a pair in such a way
that a net magnetic moment is generated that is opposite in direction to the applied
field H. The magnitude of bulk susceptibility for most of the diamagnetic materials
does not change with temperature neither it depends on the strength of the external
magnetising field H. Since application of an external magnetic field H effects the
spin and orbital motions of electrons of any atom, diamagnetic effect is produced in
all types of atoms irrespective of their class of magnetism; atoms of paramagnetic,
ferromagnetic, antiferromagnetic and ferrimagnetic materials also show diamagnetic
behaviour, but the magnitude of diamagnetic affect is much too small in these atoms
as compared to other more dominant effects. As such it may not be wrong to say that
diamagnetism is universal and more fundamental than any other type of magnetism.
3.7 Classification of Magnetic Materials 155
3.7.1.1 Langevin’s Theory of Diamagnetism
Though diamagnetism is a quantum phenomenon, a classical theory for diamag-

netism was proposed by Langevin in 1905. The classical theory was able to explain
most of the observed characteristics of diamagnetic materials. Brief outline of the
theory is presented here.
Langevin’s theory is based on the assumption that electrons circulating (around the
nucleus) in closed circular paths in an atom produce magnetic dipole pole moments,
the sum of dipole moments in a diamagnetic atom which is not subjected to any
external magnetic field is zero because of the equal number of electrons circulating in
clockwise and anticlockwise directions. However, if atoms of a diamagnetic material
are put in an external magnetic field H, they develop magnetic induction B = μH =
μ0 (H + M). Induction B applies a torque τ on each electron of the atom. Torque τ
makes each electron of the atom to undergo a rotational motion about the direction
of B. This rotational motion around B is called Larmor precession, and the angular
frequency ω L of Larmor precession is given by;
eB
ωL = (3.28)
2m e
Here e and me are respectively the charge and mass of the electron.
The angular momentum L p associated with precessional motion is given as;
⟨ ⟩
L p = m e ωL r 2 (3.29)
⟨ ⟩
where r 2 is the mean square distance from an axis through the nucleus parallel to
B.
Larmor precession motion of electron will associate an additional magnetic dipole
moment μep with each electron of the atom, given by;
e
μep = − Lp (3.30)
2m e
If each atom has Z electrons and number of atoms per unit volume (number density
of atoms) in the material is N, then the total additional dipole moment per unit volume
or magnetisation M that will develop in the material due to Larmor precession will
be given as;
( )
e e ( ⟨ ⟩)
M = Z N μep = −N Z Lp = −N Z m e ωL r 2 (3.31)
2m e 2m e
Substituting the value of ω L from Eq. (3.28) in Eq. (3.31) one gets,
( )
e eB (⟨ 2 ⟩)
M = −N Z (m e ) r (3.32)
2m e 2m e
Fig. 3.7 a Electron undergoing Larmor precession. b H–M graph for a diamagnetic material
However, B = μ0 H putting this in Eq. (3.31) gives
M e2 (⟨ 2 ⟩)
= χ = −μ0 N Z r (3.33)
H 4m e
(⟨ ⟩)
The value of r 2 may be calculated with reference to Fig. 3.7a in terms of the
mean square radius ρ of the orbit as,
ρ2 = x 2 + y2 + z2
But
ρ2
x 2 = y2 = z2 =
3
and
2 2
r 2 = x 2 + y2 = ρ
3
(⟨ ⟩) ⟨ ⟩
Therefore, r 2 = 23 ρ 2 substituting this in Eq. (3.33) gives,
( )
M e2 2 ⟨ 2 ⟩ e2 ⟨ 2 ⟩
= χ = −μ0 N Z ρ = −μ0 N Z ρ (3.34)
H 4m e 3 6m e
Expression (3.34) for the bulk magnetic susceptibility of diamagnetic materials

tells:
1. Atoms of diamagnetic materials as a whole does not possess any permanent
dipole moment
2. Magnetisation in diamagnetic materials is produced by the influence of the

external magnetic field which makes the electrons of the atom to undergo Larmor
precessional motion.
3. Diamagnetic
⟨ ⟩ materials have a small negative value of susceptibility, small
because ρ 2 is small.
4. Negative value of susceptibility means that the magnetisation produced by the
magnetising field is opposite in direction to the magnetising field. B < μ0 H .
5. Bulk susceptibility of diamagnetic materials depends on the number of electrons
per atom (Z), number density (N) of atoms (atoms per unit volume) and is inde-
pendent of temperature of the specimen and the intensity of the magnetising field
H.
6. Figure 3.7b shows the variation of M with H for a diamagnetic material.
Same results regarding diamagnetism are obtained from a quantum mechanical
approach using perturbation theory.
Some important diamagnetic elements are: H, Be, B, C, N, F, Ne, Si, P, S, Cl, Ar,
Cu, Zn, Ga, As, Se, Br, Kr, Ag, Cd, Sb, Hg, Pb, Bi, etc.
Some diamagnetic ions: A given ion is diamagnetic or not depends on its elec-
tronic configuration. If the highest electron energy states have even number of elec-
trons then the ion will be diamagnetic; for example Ca2+ ion has electron structure
1s2 2s2 2p6 3s2 3p6 has 6-electrons in the highest energy level 3p, hence it is diamag-
netic. Thus electronic configurations of atoms and ions may reveal their magnetic
properties.
Diamagnetic materials, because of the negative sign of their susceptibility, get
repelled by external magnetic fields and try to move towards the weaker part of the
external field. Diamagnetism is observed in water, wood, and in most of organic
molecules including most of living organism. The fact that water is diamagnetic may
be demonstrated by rapping a very strong permanent magnet (super magnet made
from rare earth compounds) with a very thin layer of water. The strong magnetic
field repels water which produces a small dimple (bulge) in water layer. The bulge
is small and may be observed only by a reflection microscope.
Superconductors might be considered as perfect diamagnetic material (χ = −1)
since they repel (or got repelled by) magnetic fields. Strong repulsion of strong
diamagnetic materials by super magnets may give rise to stable levitation. Pyrolytic
graphite an unusually strong diamagnetic material (χ = −40.9) can be levitated in
a stable equilibrium in the field of a permanent super magnet made from compounds
of rare earths. Levitation of living organisms, like frog, mice, etc. in the fields of
superconducting magnets has also been demonstrated.
SAQ: What causes the diamagnetic behaviour of some atoms?
3.7.2 Paramagnetic Materials
Except the atoms/ions or molecules of diamagnetic materials, atoms (ions/molecules)

of all other materials that may exhibit paramagnetism, ferromagnetism, ferrimag-
netisms or antiferromagnetism possess an inherent magnetic dipole moment; each
atom/ion is like a tiny bar magnet. The inherent magnetic dipole moment comes
essentially from the spin motion of electrons (as spin motion is twice as effective as
orbital motion in producing magnetic effects) of unpaired electrons in outer shells
of the atom/ion or molecule. Normally at temperature T > 0 K the inherent magnetic
dipoles have random orientations because of thermal agitation. The resultant magneti-
sation, which is the vector sum of atomic dipole moments, is zero because of the
random distribution of atomic dipoles and also for minimising the magnetic energy
of the system. However, when some external magnetic field H is applied, randomly
oriented inherent dipoles try to align in the direction of the applied magnetic field. At
a given temperature T, each atomic dipole experience two torques: a torque τ H due
to field H that tries to align the atomic dipole in the direction of H and a torque due
to thermal motion τT that opposes the torque τ H of alignment. Since torque τ H has
different magnitudes for different atomic dipoles all atoms of the specimen do not
align completely, resultant magnetisation M generated from the partial alignment of
inherent dipoles has a small value but it is in the same direction as that of H. This
results in a small positive value for the bulk magnetic susceptibility χ . Materials
with positive small values of bulk susceptibilities in the range of + 10–5 to + 10–3
are termed as paramagnetic. Large number of elements like O, Na, Rb, Sc, Al, Sn,
etc. is paramagnetic.
When atoms of a paramagnetic material are put in an external magnetic field H,
two happenings take place; (i) the inherent magnetic dipole moments of atoms try to
align along the direction of the field H giving rise to a magnetisation, say M 1 in the
direction of H and (ii) the electrons of the atom starts Larmor precession around the
direction of induction B which generates a magnetisation − M 2 in a direction opposite
to H. Since M 1 is generally much larger than M 2 , the net result is a magnetisation of
magnitude + (M 1 − M 2 ) pointing in the direction of H. Obviously, M 2 is a reflection
of the diamagnetic behaviour the paramagnetic material.
3.7.2.1 Langevin’s Theory for Paramagnetism
Different theories have been proposed to explain paramagnetism of different types

of materials. Langevin’s theory is applicable to those materials in which the inherent
magnetic dipoles of neighbouring atoms do not interact with each other. This is often
referred as theory for non-interacting localised electrons and holds for materials like
hydrated salts of transition metals like CuSO4 .5H2 O and gases. The diamagnetic
water molecules in such hydrated salts shield one magnetic dipole from the other
while in dilute gases atomic magnetic dipoles are far apart to be in the interaction
range.
Fig. 3.8 Partially aligned

inherent atomic magnetic
dipole
Langevin assumed that each atom/ion of the material has a permanent or inherent
dipole moment μ and that number density of atoms/ions in the specimen is N per
unit volume. Further, it is assumed that there is no interaction between these inherent
dipoles. On application of an external magnetic field H, each inherent magnetic
dipole experiences a torque τ which tries to align the magnetic dipole in the direc-
tion of external magnetic field. However, on account of the thermal motion (due
to temperature T of the specimen), all atomic dipoles could not align completely
along the direction of H. Different atomic dipoles will align to different degrees with
respect to the field H. Let us assume that a typical inherent dipole align itself at an
angle θ with respect to the direction of the external magnetic field H (see Fig. 3.8).
The potential energy E θ of this magnetic dipole in magnetic field H is given as,
E θ = μ.H = −μ.H cos θ (3.35)
The number n of atomic magnetic dipoles with energy E θ at temperature T of the

specimen may be calculated using Maxwell Boltzmann statistics which gives,
( ) ( ) ( )
− −μH cos θ μH cos θ
E
− k θT
n = n0e β = n0e k T
β = n0e kβ T
(3.36)
Here n0 is a constant and k β is Boltzmann constant.

The number dn of magnetic dipoles that have their inclinations between θ and (θ
+ dθ ) with respect to the direction of H may be obtained by differentiating expression
(3.36) with respect to θ. Therefore,
( ) ( μH cos θ )
μH
dn = n 0 e kβ T sin θ dθ (3.37)
kβ T
Each of the partially aligned magnetic dipole contributes an additional magnetic

moment m = μ cos θ to the specimen in the direction of the applied field H (compo-
nents of μ in direction perpendicular to the field H from different dipoles will cancel
each other). The average value of the total additional magnetic moment, denoted
by ⟨m⟩ may be obtained by dividing the sum of the components in field direction
contributed by each atomic dipole by the total number of dipoles,
[ ( ) ( μ cos θ ) ]
{π
{π μ cos θ n 0 kμH e kβ T
sin θ dθ
μ cos θ dn 0 βT
⟨m⟩ = 0
{π = ( ) ( μ cos θ )
{π
0 dn n0 μH
e kβ T sin θ dθ
0 kβ T
Or
[ ( μH cos θ ) ]
{π
0 μH cos θ e kβ T sin θ dθ
⟨m⟩ = ( ) (3.38)
{π μH cos θ
0 e kβ T
sin θ dθ
Let us make following substitutions in Eq. (3.38);
μH
y= and x = cos θ ; then dx = − sin θ dθ
kβ T
Now, when θ = 0, x = 1 and when θ = π, x = −1.

With above substitutions and change of limits, expression (3.38) becomes,
[ ]1
{1 μ xy e yx − y12 e yx
μ yx
−1 xe dx −1
⟨m⟩ = {1 = [ ]1
−1 e yx dx 1 yx
e
y −1
Or
[ y ]
−y y
μ ey + e y − ey 2 + e−y [⎧ y ⎫ ]
y2 e + e−y 1
⟨m⟩ = [ y ] =μ −
e −y
− ey e y − e−y y
y
Or
( )
1
⟨m⟩ = μ coth y − = μL(y) (3.39)
y
Here L(y) is called Langevin function.

The average magnetic moment ⟨m⟩ multiplied by the number of atoms per unit
volume N (number density of atoms in the specimen) gives the additional magneti-
sation M generated due to partial alignment of inherent magnetic dipoles under the
influence of applied external magnetic field H. Hence;
M = ⟨m⟩N = N μL(y) (3.40)
Depending on the value of parameter y there may be two possible situations;

(i) For y ≫ 1, i.e. for large values of y
Y ≫ 1 means that (y =) kμH

βT
≫ 1 or μH ≫ kβ T
The above condition corresponds to the case of high value of applied magnetic
field H and very low temperature T of the specimen, the function L(y)
approaches 1 and M becomes.
M Nμ
M = N μ and susceptibility χ = = (3.41)
H H
It may be noted that in the case μH ≫ kβ T susceptibility of a paramagnetic

substance is independent of temperature T of the specimen and is given by
Nμ
H
. The value of magnetisation (M = Nμ) of a paramagnetic substance at
low temperature and high magnetising field H is called the saturation value of
magnetisation, often denoted by M s and corresponds to the complete alignment
of all the inherent dipole moments of the specimen along the direction of the
applied field H.
When y ≪ 1, i.e. μH ≪ kβ T
For y ≪ 1, e y and e−y may be expanded in the form of converging series as
fellows,
y2 Y3
ey = 1 + y + + + ···
2! 3!
and
y2 Y3
e−y = 1 − y + − + ···
2! 3!
With these substitutions, Langevin function L(y) becomes

⎡ ( ) ⎤
[⎧ −y ⎫ ] 2 1+ y2
+ y4
···
e +e
y
1 2! 4! 1⎦
L(y) = − =⎣ ( )−
e y − e−y y 2 y+ y3
+ y5
... y
3! 5!
Since successive terms of the two series appearing in the numerator and the
denominator of the above expression decrease very fast, it is enough to retain
first two terms of the series to obtain the approximate value of function L(y).
Or
[ ] ⎡( 2
) ⎤
1 + y2
2
1 1 + y2 1
L(y) ∼
= y3
− =⎣ ( )− ⎦
y + 3.2.1 y y
2
y 1 + y6
[( )( )( )−1 ]
1 y2 y2 1
= 1+ 1+ −
y 2 6 y
Or
[( )( )(( )) ]
1 y2 y2 1
L(y) ∼
= 1+ 1− −
y 2 6 y
( )[ 2 2 4 ]
1 y y y
= 1+ − − −1
y 2 6 12
Dropping the y4 term one gets,

( )( 2 )
1 y y
L(y) ∼
= =
y 3 3
Thus magnetisation per unit volume
N μy μH
M = N μL(y) = = Nμ
3 3kβ T
and magnetic susceptibility
M N μ2 C
χ= = = (3.42)
H 3kβ T T
N μ2
where C = 3kβ
is called Curie constant.
Expression (3.42) is called Curie law and tells that for the condition μH ≪ kβ T
the magnetic susceptibility of paramagnetic materials is inversely proportional to the
temperature.
Variation of the ratio of magnetisation M to the saturation magnetisation M s with
the value of the parameter y (= μH/k β T ) is shown in Fig. 3.9. It may be noted that
for low temperature and high magnetising field, the ratio M/M s increases with H as
the tangent to the curve for low H and high temperature.
As already mentioned, the above derivation of susceptibility for paramagnetic
materials is applicable only to those atoms/ions or molecules where there is no inter-
action between nearby inherent magnetic moments of the material. However, this
assumption does not hold, particularly, for paramagnetic metals which contain large
number of free electrons in conduction band. Application of an external magnetic
field affects both the magnetic moments due to the orbital and spin motions in such
cases. Since spin of an electron is totally a quantum mechanical concept, Pauli model
of paramagnetism is a quantum mechanical model that is beyond the scope of this
discussion. In Pauli model conduction electrons are considered essentially free and
Fig. 3.9 Variation of the

ratio M/M s with parameter y
under the applied external magnetic field H an imbalance between electrons with
opposite spin is setup that leads to a low value of magnetisation in the direction of
field H. This imbalance may be understood in terms of Fig. 3.10. Figure 3.10a shows
the Fermi–Dirac distribution of electrons with opposite spins before the application
of magnetising field H. As may be observed in this figure the two halves of the distri-
bution for opposite spins are equal. However, on application of external magnetic
field H, the energy of component with spin parallel to H decreases by the amount
μb .H , while those with spin opposite to H increases by the same amount. Here
μb represents the inherent magnetic dipole moment of the atom/ion. As a result of
pulling down (in energy) of parallel spin distribution, some electrons from antipar-
allel side in the neighbourhood of Fermi level, fall into parallel spin side, increasing
the number of electrons with parallel spin over the number of electrons with antipar-
allel spins. This slight imbalance in the number of electrons on the two sides gives
rise to a small value of magnetisation in the direction of the applied field.
The susceptibility, though essentially independent of the temperature but in case
of paramagnetic metals, may have some temperature dependence because of the
change in electronic band structure due to field H.
Some metals like Al, Mg, Ti, V, etc.; some diatomic gases like O2 , and NO,
ions of transition metals and rear earth metals and their salts along with rear earth
oxides show paramagnetism. Many minerals and other materials found in nature are
paramagnetic, for example pyrrhotite (Fe3 S8 ), ilmenite (FeTiO3 ), siderite (FeCO3 ),
quartz (SiO2 ), etc. show paramagnetic behaviour.
SAQ: Use expression (3.36) to calculate the number of dipoles that will have zero
energy in the applied field H.
Fig. 3.10 Fermi–Dirac distribution of electrons with spin up and down a before the application of
magnetic field H, b after the application of magnetic field H
3.7.3 Ferromagnetic Materials
Ferromagnetic materials are characterised by a large positive value of susceptibility

(0.1–5 × 103 ) below a certain temperature, called Curie temperature or transition
temperature or critical temperature. Ferromagnetic materials are made up of atoms/
ions that have inherent permanent magnetic dipole moments (like paramagnetic
atoms essentially because of the spin of unpaired electron in valence orbital) and
are arranged in a lattice such that dipole moments of atoms/ions in a group align
parallel to each other. Ferromagnetism is the outcome of some sort of ordering of
atomic/ionic magnetic dipoles in different parts of the material.
3.7.3.1 Weiss’s Theory of Ferromagnetism
French physicist Pierre Ernest Weiss in 1907 proposed a phenomenological theory

to explain ferromagnetism. The theory is based on the assumption that neighbouring
atomic magnetic dipoles, due to certain mutual exchange interactions, align them-
selves in groups, forming several very small regions of volume, called domains.
Atomic (ionic or molecular) dipoles in a domain are all aligned parallel to each other
generating a magnetisation Ms for each domain. Since the number of atomic dipoles
and their direction of alignment are different for different domains, the magnitude and
direction of magnetisation Ms has a different value for each domain. The resultant
magnetisation M R of a given specimen of a ferromagnetic material may be obtained
by taking vector sum of magnetisation of individual domains;
Σ
MR = Ms
all domains
At a temperature below Curie temperature T c , and in absence of any external

magnetic field, magnetisations of different domains are randomly oriented so that
the overall magnetisation M R is zero and the material does not show any magnetism.
This is required for minimisation of magnetic energy of the specimen in absence of
any external magnetic field.
Although there may be large number of domains in a piece of a ferromagnetic
material but for simplicity, the total volume of the specimen may be divided into
four equal volumes having dipole moments aligned along the four basic directions
as shown in Fig. 3.11c. The directions of magnetisation in different domains form
a closed loop to minimise magnetic energy of the system and to avoid the leakage
of magnetic flux, as shown in this figure. On application of an external magnetising
field H, domains (and individual atomic dipoles) rotate in an effort to align their
magnetisation in the direction of applied field H. Out of the large number of domains
some domains already have their magnetisations in the direction of applied field H,
such domains are called favourable domains. When magnetising field H is switched
on the size (volume) of favourable domains in the specimen increases at the cost of
the volume of unfavourable domains as shown in Fig. 3.11d. Increase in the size of
favourable domains is achieved by the motion of domain walls.
This re-alignment of domains results in generating a net magnetisation in the
direction of field H. When temperature is increased, the thermal motions of atoms
destroy domain structure and randomise the direction of atomic dipoles, resulting in
zero magnetisation.
The characteristic feature of the ferromagnetic order in each domain is a sponta-
neous magnetisation M s due to spontaneous alignment of atomic magnetic moments.
The spontaneous magnetisation M s tends to lie in a specific direction determined by
the shape and/or crystal structure. Further, the spontaneous magnetisation disappears
beyond a certain temperature called Curie temperature. Weiss theory assumes that
the spontaneous magnetisation M s of each domain is due to an internal ‘molecular’
magnetic field H i . Which is proportional to the spontaneous magnetisation of the
domain, i.e.
H i = nw M s (3.43)
The constant of proportionality n w , called Weiss molecular field constant, is a

measure of the extent of alignment of atomic dipoles in a domain. In order to
completely align all atoms of a domain in a particular direction (n w = 100%),
the internal molecular field H i must be quite large. The origin of such a large
internal molecular field remained a mystery until Heisenberg introduced the idea
of the exchange interactions in 1928.
If H is the external magnetising field then the effective magnetic field acting on
each atom or ion may be written as
H e f f = H + H i = H + nw M s (3.44)
Fig. 3.11 Magnetic domains in a ferromagnetic material (a) in absence of any external field (b) in
presence of external magnetic field H. c In absence of any magnetising field area/volume occupied
by domains in four basic directions is same. d On applying the magnetising field H, the area/volume
of the domain having magnetisation parallel to the applied field H increases
Let N be the number of atoms per unit volume, J the total angular momentum
quantum number of each atom, then the possible components of magnetic moment
is,
M j gμB
where M j may have values, M j = J, ( J − 1), (J − 2) . . .−(J − 2), −(J − 1), −J

and g is Lande’s splitting factor while μB is Bohr magneton.
Hence, the potential energy of a dipole with component M j gμ B along H is
P E = −M j gμ B .H (3.45)
Now from statistical mechanics it follows that the total magnetic moment per unit
volume or magnetisation M along H is given by
(M )
√+J j gμB
kβ T
N −J M j gμB e
M= (M ) (3.46)
√+J j gμB
kβ T
−J e
where kβ is Boltzmann constant.

Equation (3.46) may also be written as;
M = N g J μB B J (x) (3.47)
where
( ) ( x )
2J + 1 2J + 1 1
B J (x) = coth x− coth (3.48)
2J 2J 2J 2J
and
g J μB Heff g J μB (H + n w M)
x= = . (3.49)
kβ T kβ T
In case when there is no external magnetic field (H = 0), H eff = nw Ms and

Eq. (3.49) reduces to
g J μB (n w Ms )
x= (3.50)
kβ T
Or
kβ T x
Ms (T ) = (3.51)
g J μB n w
When temperature T goes to zero, T → 0, x → ∞ and B J (x) → 1 all atomic

dipole magnetic moments of a domain completely align themselves parallel to the
direction of magnetisation M s .
It follows from Eq. (3.47) [M = N g J μB B J (x)], that in case when T →
0, B J (x) → 1, therefore,
Ms (0) = N g J μ B (3.52)
Dividing Eq. (3.51) by Eq. (3.52) one gets,
Ms (T ) kβ T x
= ( )2 (3.53)
Ms (0) N n w g J μβ
Also, when one divides Eq. (3.47) by Eq. (3.52), one gets,
Ms (T )
= B J (x) (3.54)
Ms (0)
The value of M s (T ), spontaneous magnetisation at temperature T, may be found

by solving Eqs. (3.53) and (3.54) simultaneously. Simultaneous solution to these
equations may be obtained in two different ways; either by graphical method or by
numerical analysis. Let us find the solution using graphical method.
Figure 3.12 shows the plots of Eq. (3.53) for four values of temperature T; T >
T C , T = T C , T 1 < T C and T 2 < T C . The function BJ (x) has also been plotted for J
= 1/2 value in the figure. As may be observed in the figure, a solution of the two
equations that gives a real value greater than zero for the ratio M s (T )/M s (0) occurs
only for T < T C , when the two curves cut each other. For T > T c and T = T c the two
curves cut each other only at t = 0. It means that a positive value of the ratio M s (T )/
M s (0) may be obtained only when the temperature is below the transition temperature
or Curie temperature T C . Moreover, for temperatures less than T c , the ratio M s (T )/
M s (0) decreases as temperature increases, as one goes from T 2 to T 1 . It may be noted
that in Fig. 3.12 temperature increases as one moves from right to left. Hence, for
any temperature T < T c , the ratio M s (T )/M s (0) is inversely proportional of temper-
ature, i.e. the spontaneous magnetisation decreases with increasing temperature and
vanishes beyond temperature T C called ferromagnetic Curie temperature.
Variation of the ratio M S (T )/M S (0) with temperature ratio (T /T C ) is shown in
Fig. 3.13. The maximum value of spontaneous magnetisation occurs at absolute zero
temperature where there is no thermal motion of atoms and all atomic dipoles in a
given domain line-up in one direction under the influence of internal exchange inter-
action field Hi . Spontaneous magnetisation decreases with the increase of tempera-
ture becoming zero at Curie temperature T C . Since M S (T )/M S (0) depends on function
Fig. 3.12 Graphical solution

of Eqs. (3.53) and (3.54) for
J = 1/2 to find the
spontaneous magnetisation
M s when T < T C .
Equation (3.53) is also
plotted for T > T c and T =
TC
BJ (x) from Eq. (3.54), the ratio depends on the total momentum J associated with
the atom and hence rate of fall for the ratio is J dependent as shown in Fig. 3.13.
In case when T > T C , when there is no spontaneous magnetisation, the material
behaves as a paramagnetic substance and a small external field may be required to
align some of the atomic dipoles to produce some magnetisation. The external field
must be small to avoid the state of saturation. Now from Eq. (3.49)
g J μB Heff g J μB (H + n w M)
x= = ≪ 1; since T is large and H is small
kβ T kβ T
Further, from Eq. (3.47)
M = N g J μB B J (x)
( J +1 )
But for small x ≪ 1, B j (x) ∼
= 3J
x; putting this value in the expression above,
one gets;
( ) ( )
J +1 J + 1 g J μB (H + n w M)
M = N g J μB B J (x) ∼
= N g J μB x = N g J μB
3J 3J kβ T
Or
N g 2 μ2B ( J + 1)
M= [H + n w M] (3.55)
3kβ T
But magnetic susceptibility χ = M

H
, therefore dividing Eq. (3.55) by H gives;
[ ]
M N g 2 μ2B (J + 1) M N g 2 μ2B (J + 1)
χ= = 1 + nw = [1 + n w χ]
H 3kβ T H 3kβ T
Or
Fig. 3.13 Variation of ratio

M S (T )/M S (0) with
temperature ratio T /T c
( )
N g 2 μ2B (J + 1)n w N g 2 μ2B (J + 1)
1− χ= (3.56)
3kβ T 3kβ T
Defining
N g 2 μ2B (J + 1)n w
Tc ≡ (3.57)
3kβ
Substituting the above value of T c in Eq. (3.56) gives;

( ) ( )
TC Tc TC 1 C
1− χ= Or χ = =
T nw T nw T − Tc T − Tc
Or
C
χ= (3.58)
T − Tc
( )
Here C = nTCw is a constant for a give material and is called Curie constant.
Equation (3.58) defines Curie–Weiss law and gives the magnetic susceptibility
of a ferromagnetic material above Curie temperature for low external fields. The
variation of magnetic susceptibility of materials at temperatures higher than Curie
temperature, where a ferromagnetic material behaves as a paramagnetic substance, is
well explained by Curie–Weiss law. The Curie temperature Tc for different material
may be determined experimentally by measuring susceptibility at different temper-
atures. Once T C is known other parameters of the material may also be determined;
for example in case of ferromagnetic element Gd, the experimentally determined
value of T C is 292 K; J = S = 7/2, g = 2, N = 3.0 × 1028 m−3 . This data gives
M s (0) = Ng μB J = 1.95 M A m−1 and Bi = μ0 H i = 144 T.
Figure 3.14 shows the temperature dependence of susceptibility for a paramag-
netic material and for a ferromagnetic material. A ferromagnetic material undergoes
phase transition at Curie temperature T C . This results in singularities in the behaviour
of physical properties like susceptibility, magnetisation, specific heat, etc.
(i) Exchange interactions
It can be shown that interactions between atomic magnetic dipoles cannot
generate a magnetic field strong enough to align all atoms of a domain in a
particular direction, i.e. dipole interactions are not strong enough to generate
the internal magnetic field H i ≈ 100T which is required to align atoms in
a domain. Weiss in 1907, while proposing his theory simply assumed that a
sufficiently strong internal field H i responsible for ferromagnetism is present
in each domain, without giving any explanation for its generation. The riddle
was solved in 1928 when Heisenberg introduced the concept of exchange inter-
actions. Exchange interactions originate from electrostatic Coulomb repulsion
between electrons of the neighbouring atoms and Pauli’s exclusion principle.
Fig. 3.14 Temperature dependence of a, c paramagnetic substance b, d Curie–Weiss susceptibility

of a ferromagnetic material
Pauli’s exclusion principle forbids two electrons in a given energy state to

have same values of all their quantum numbers. Since production of exchange
interactions are strictly quantum mechanical phenomenon, which is beyond the
scope of this discussion, we will not go into its details. Heisenberg derived the
following Hamiltonian, H for exchange interaction
H = −2 j S1 S2
Here S1 and S2 are the spins of neighbouring atoms, and j is the exchange
integral (do not confuse this J with total spin J). J > 0 indicates a ferromagnetic
interaction favouring parallel spin alignment (↑↑) while J < 0 indicates an
antiferromagnetic interaction favouring antiparallel spin alignment (↑↓).
SAQ: What happens to the domain walls when the magnitude of the external
field H is changed?
(ii) Spin wave
The lowest energy state of a ferromagnetic system occurs when all spins (spins
of all atoms) are parallel to each other along the direction of magnetisation.
However, when one of the spins tilts or get disturbed, it begins to precess
around the direction of magnetisation. The disturbance so produced propa-
gates as a wave through the system because of exchange interaction between
neighbouring atoms as shown in Fig. 3.15a.
Fig. 3.15 A spin wave on the line of spin
Spin waves are analogues to lattice waves created by the oscillation of atoms
about their equilibrium position. In spin waves spins precess around equilib-
rium magnetisation and precession of atoms are correlated through exchange
interactions. Spin waves may be quantised, like quantisation of lattice waves
with quanta called phonon. The quantised spin wave is called magnon.
Elements like iron (Fe), cobalt (Co), nickel (Ni), gadolinium (Gd) and dyspro-
sium (Dy) are ferromagnetic at room temperature. The Curie temperature for
Fe, Co, Ni and Gd is respectively, 770 °C, 1131 °C, 358 °C and 565 °C; however,
EuO has Curie temperature of 70 K (343 °C) and EuS even lower.
(iii) Saturation magnetisation Msat
Three parameters, namely the Curie temperature, saturation magnetisation and
magnetocrystalline anisotropy, are called intrinsic properties of a magnetic
material as they do not depend on the microstructure, i.e. on the grain size and
grain orientation in the crystal.
Saturation magnetisation (M sat ) tells about the maximum magnetic field that
may be produced by a material. Msat depends on three factors: (a) strength
of each atomic/ionic magnetic dipole (m), (b) packing density of atomic/ionic
dipoles and (c) the degree of alignment of dipoles at a given temperature.
Factor (a) depends on the nature of the atom and its electronic configuration
while factor (b) is determined by crystal structure and the presence of any non-
magnetic elements within the structure. Factor (c), the degree of alignment of
atomic/ionic dipoles depends on temperature; higher the temperature less will
be alignment, and on magnetic anisotropy of the crystal.
(iv) Magnetic anisotropy
In crystalline magnetic materials it is often observed that there is one particular
crystallographic direction magnetisation along which is easier as compared to
the other directions. For example, in case of the hexagonal crystal structure of
cobalt (Co) it is easy to magnetise in the crystal direction (001) as compared to
any other direction. It is hard to align the magnetic dipoles along the directions
⟨1010⟩ which lie in the plane normal to the crystal axis 001). See Fig. 3.16b. This
anisotropy is created by the coupling of electron orbitals to the crystal lattice. In
Fig. 3.16 Easy and hard

directions of magnetisation
in hexagonal cobalt crystal
the easy direction of magnetisation this coupling is such that electron orbitals
are in their lowest energy state.
Magnetic anisotropy of a crystal is measured in terms of the anisotropy field
Ha defined as the magnetic field needed to saturate the magnetisation in a hard
direction.
(v) Magnetic hysteresis in ferromagnetic substances
A ferromagnetic material in absence of any external magnetic field behaves
as a diamagnetic material because of random orientation of magnetisations of
different domains. However, application of a small external magnetic field H a
large magnetic induction B or magnetisation M in the direction of the applied
field H is produced. When a ferromagnetic material is magnetised in a particular
direction, it does not revert back to the state of zero magnetisation when the
external magnetising field H is withdrawn. In order to bring a magnetised
ferromagnetic specimen back to the state of zero magnetism, a magnetic field
opposite in direction to the field H has to be applied.
Figure 3.17 shows the behaviour of a ferromagnetic specimen below its Curie
temperature, subjected to an alternating magnetic field. The strength of the
magnetising field H is plotted on the X-axis while the magnitude of the magnetic
flux density B or the strength of magnetisation M (B = μ(H + M)) is shown
on the Y-axis. The starting point is the origin O when H = 0 and B = 0. As
H is increased in a given direction the flux density B in the direction of H
also increases reaching the point a (H 0 , B0 or M 0 ). Any further increase of H
beyond H 0 does not increase B or M above B0 (or M 0 ). This is called the state
of saturation. In the state of saturation, magnetisations of all domains in the
specimen get maximum aligned in the direction of field H. When magnetising

field H is reduced from the saturation value H 0 , the flux density B also reduces
but it does not follow the path traced while increasing H. The variation of
B with reducing H is shown by the part of the curve shown by abcd in the
figure. The important aspect of the curve is that when H = 0 at point b, there
remains a residual magnetisation Br or M r in the specimen. This indicates that
the ferromagnetic specimen retains some sort of a memory about its previous
magnetic history which persists even when the magnetising field is withdrawn.
This happens because thermal motion of atoms does not fully randomise the
magnetisation of magnetic domains (which were aligned in the direction of H)
in the specimen and they remain partially aligned in the direction of field H
even when H becomes zero. Now if the external magnetising field is applied in
opposite direction (− H), the magnetic flux density B or magnetisation M of the
specimen decreases from the value Br (or M r ) to a zero value at point C on the
graph where the value of magnetising field is (− H c ). This value (− H c ) is called
the coercive field value that ultimately removes any residual magnetisation of
the specimen. Further increase in the magnitude of magnetising field in opposite
direction (− H) generates magnetisation in the specimen in the direction along
(− H) which reaches the saturation value (− B0 or− M 0 ) in opposite direction
at point d on the graph. When the magnitude of magnetising field (− H) is
reduced, B (or M) in the specimen follow the path indicated by the part defa of
the curve, where point e corresponds to residual magnetisation (− Br or − M r )
in opposite direction and point f to coercive field (H c ) that is required to nullify
any residual magnetisation of the specimen. On increasing H beyond H c , the
magnetisation of the specimen increases and ultimately attains the saturation
value at point a. The closed curve abcdefa shown in the figure is often called the
B–H curve or hysteresis loop of the ferromagnetic specimen. The area enclosed
by the hysteresis loop gives the amount of energy consumed in one cycle of
magnetisation from zero to saturation value in one direction to the saturation
vale in the other direction and back. The energy consumed in magnetisation per
cycle given by the area of the B–H curve generally appears in the form of heat.
This may be seen in case of a solenoid that is fed by alternating (ac) current
generating magnetic field with alternating polarities at its axis, a ferromagnetic
material like iron rod if placed at the core of the solenoid gets heated. Similarly,
magnetic cores of transformers that undergo repeated cycles of magnetisation
get heated.
The shape, particularly, the area enclosed by the hysteresis loop has different
values for different types of ferromagnetic materials. Ferromagnetic materials
with large area of B–H loop are called Hard magnetic materials; they have
large value of residual magnetism, large value of coercive field and consume
large amount of energy in each cycle of magnetisation. Natural iron is such
a material; hard magnetic materials are used to make permanent magnets on
account of the large value of residual magnetisation and large value of coer-
cive field. On the other hand materials with narrow hysteresis loop, that have
Fig. 3.17 Magnetic hysteresis loop for a ferromagnetic material at T < T C
low value of residual magnetisation and coercive field are called soft magnetic
materials. Soft magnetic materials may be easily magnetised and demagnetised
without loss of much energy. Soft materials have high value of relative perme-
ability (μr ) and low value of coercive field. Silicon steel is a very good example
of soft magnetic materials and is often used as core material in solenoids, trans-
formers, and relays that operate on alternating current generating alternating
magnetic fields.
Common soft magnetic materials are iron, iron–silicon alloys (with 1–5%
silicon) and nickel–iron alloys, also called permeability alloys, with preferred
nickel contents of 42–79%. Addition of molybdenum gives extra electrical
resistivity and addition of copper results in higher initial permeability. Soft
magnetic ceramics, also called ceramic magnets, have been originally made
from iron oxide (Fe2 O3 ) with one or more divalent oxides like that of ZnO,
MgO or NO. The mixture of these oxides is first calcined and grinded to powder,
pressed to the desired shape and sintered. Vectolite is typical light weight and
very high resistivity (like that of an insulator) magnet made by moulding ferric
and ferrous oxides and cobalt oxide. Magnadur (BaO.Fe2 O3 ), made from
BaCO3 (barium carbonate) and ferric oxide is also a soft magnet material.
Table 3.2 lists the properties of some important soft magnetic materials.
Magnets made from hard magnetic materials have strong resistance against
demagnetisation (large coercive field) and large area of hysteresis loop. Details
of some important hard magnetic materials are tabulated in Table 3.3.
Table 3.2 Properties of some soft magnetic materials

Material/trade % Composition by weight, Relative permeability Coercive field
name remainder iron initial maximum (Oe)
High-purity iron Impurity < 0.05% 10,000 200,000 1.1
Commercial iron Impurity 0.2% 250 9000 0.9
Transformer, 2.2% Si 1500 7000 0.35
M-15
Armature M-43 0.95% Si 100 4100 0.94
Table 3.3 Hard magnetic materials

Material Composition Curie Coercive Residual Preparation Mechanical
(% by weight) temperature field H c induction properties
(°C) (Oe) Br (T)
Alnico-5 50Fe, 24CO, 900 620 1.25 Casting and Hard and
15Ni, 8Al, annealing brittle
3Cu
Alnico-8 34Fe, 35Co, 860 1600 0.83 Hard
15 Ni, 7Al,
5Ti, 4Cu
Mn–Al–C 70Mn, 29Al, 300 2700 0.61 Casting, Hard
0.5Ni, 0.5C extruding,
annealing
Ba ferrite Bao.6Fe2 O3 450 2100 0.43 Pressing, Brittle
sintering
3.7.4 Antiferromagnetic and Ferrimagnetic Materials
Fifteen elements of the periodic table, namely O, Cr, Mn, Fe, Co, Ni, Nd, Sm,
Eu, Gd, Tb, Dy, Ho, Er and Tm in their solid states show some sort of magnetic
order. As expected, atoms of all these elements have unpaired electrons and associ-
ated magnetic dipole moments essentially from the spins of the unpaired electrons.
Heisenberg has shown that the magnetic order in these fifteen elements originates
from the exchange interactions between the electron clouds of atoms/ions/molecules
subject to Pauli’s exclusion principle (i.e. from electrostatic interactions) and not
from the mutual interaction between the magnetic dipoles of atoms or from the spin–
orbit interactions in atoms, etc. The electrostatic exchange interactions may align
the spin magnetic dipole moments of individual atom/ion/molecule either parallel
or antiparallel to each other. For example, in case of H2 molecule the energy of
the parallel (spin) alignment of two atoms (↑↑), called triplet state, and antiparallel
alignment (↑↓), called singlet state, have different value that depends on the rela-
tive separation of the two atoms as shown in Fig. 3.18. Therefore, the singlet or the
triplet alignment in the ground state of a martial depends on their spins and relative
separation of atoms/ions, etc. in the crystalline structure.
Fig. 3.18 Singlet and triplet

state energies for H 2
molecule
According to Heisenberg theory of electrostatic exchange interactions the

effective interaction between a pair of atoms/ions of spins S 1 and S 2 is given by
⇀ ⇀ 1 3
S1 . S2 = for triplet; and − for singlet,
4 4
And the interaction energy U is given as
( )⇀ ⇀ 1( )
U = − E singlet − E triplet S1 . S2 + E singlet + 3E triplet
4
Or
⇀ ⇀
U = −J S1 . S2 + Constant (3.59)
Here E singlet and E triplet respectively, represents the energies of the singlet and triplet
states. Factor J in Eq. (3.59), called ‘exchange coupling constant’ may have positive
(J > 0) or negative (J < 0) values. J > 0 refers to the case when spin will orient in the
same direction (triplet case) while J < 0 refers to the antiparallel alignment of spins
(singlet case). In ferromagnetic materials spins and magnetic dipoles of adjacent
atoms/ions in a given domain align parallel to each other, and it refers to the case
when J > 0. In antiferromagnetic and ferrimagnetic materials, J < 0 and spinful atoms/
ions of such materials align their spins/magnetic moment in opposite directions in
a domain. It may therefore be said that both ferromagnetic and antiferromagnetic/
ferrimagnetic materials have domains, which are created because of the electrostatic
exchange interactions but in ferromagnetic materials the spin or the magnetic dipole
moments of neighbouring atoms/ions are aligned parallel to each other while in
antiferromagnetic and ferrimagnetic materials the dipole moments of adjacent atoms/
ions are aligned opposite to each other in each domain.
Antiferromagnetism was predicted by French scientist Louis Neel in 1936. It
is experimentally difficult to detect antiferromagnetic material because these mate-
rials above a certain temperature (called Neel temperature) behave like a paramag-
netic material above Curie temperature; and show no magnetism. The final direct
confirmation of Neel’s theory was done by Harry Shull using neutron diffraction in
1938.
Figure 3.19 shows the spin (and magnetic moment) alignments of different types of
magnetic materials in their solid states. Figure 3.19a shows that atoms/ions/molecules
of the material have no net spin and magnetic moment, since they have no unpaired
electrons and, therefore, in absence of an external magnetising field they show no
magnetism. However, when an external magnetising field H is applied, the material
develops weak magnetisation opposite to the direction of the applied field H. They
are diamagnetic. Susceptibility of diamagnetic substance has a negative small value
and generally independent of the temperature. Figure 3.19b shows a paramagnetic
material; atoms/ions/molecules of a paramagnetic material have unpaired electrons
that give rise to an inherent magnetic dipole moment to each of them. Individual
dipoles do not interact with each other. At room temperature (T > 0 K) and in
absence of magnetising field H, the atomic/ionic/molecular dipoles are randomly
distributed (to minimise the magnetic energy of the system) and a paramagnetic
material does not show any magnetisation. However, application of magnetising
field H results in partial alignment of dipoles in the direction of field H generating
a magnetisation in the material. The magnitude of magnetisation increases with the
increase of H reaching a saturation value when all dipoles in the specimen get aligned
to field H. Susceptibility of paramagnetic materials is positive but small; variations
of susceptibility χ and 1/χ of paramagnetic materials with temperature are shown
in Fig. 3.20.
Atoms/ions/molecules of ferromagnetic, antiferromagnetic and ferrimagnetic
materials have unpaired electrons and, therefore, each atom/ion/molecule possesses
a dipole moment (like that of paramagnetic material), but these dipoles interact
with each other (unlike paramagnetism). Mutual interaction between dipoles arises
from electrostatic exchange interaction. The exchange interaction does two things:
(i) it creates domains in the specimen and (ii) spins or magnetic moments of all
atoms in a domain are either aligned parallel or anti-parallel. Materials for which
spins are aligned parallel are called ferromagnetic and those where spins are aligned
anti-parallel are either antiferromagnetic or ferrimagnetic.
In ferromagnetic materials all atoms/ions, etc. have their magnetic moments
aligned parallel to each other in a given domain that gives each domain a finite value
of magnetisation. In absence of magnetising field H and below Curie temperature T c ,
the orientations and magnitudes of magnetisations of different domains are random,
such that the net magnetisation of the material is zero. When external magnetising
field H is switched on, the domain walls in the sample move so as to increase the size
Fig. 3.19 Alignment of spin/magnetic moments in different types of magnetic materials
Fig. 3.20 Variation of magnetic susceptibility with temperature

of the favourable domain that results in the generation of a magnetisation in direc-

tion of H in the ferromagnetic material. Susceptibility of ferromagnetic materials
is positive and large and it decreases with temperature. Temperature dependence of
susceptibility and its reciprocal for ferromagnetic materials are shown in Fig. 3.20.
It may be mentioned that above Curie temperature a ferromagnetic material behaves
like a paramagnetic material.
Both in antiferromagnetic and ferrimagnetic materials the spin (or dipole moment)
alignment of neighbouring atoms/ions are in opposite directions as shown in
Fig. 3.19d, e. The difference between the two classes of magnetism lies in the relative
magnitudes of the magnetic moments of the adjacent atoms/ions. In case of antifer-
romagnetic materials the magnetic moments of adjacent atoms are almost equal and,
therefore, in absence any magnetising field H and below Neel temperature T N a
specimen of antiferromagnetic material does not show any magnetisation. However,
in presence of the magnetising field H below T N , domain walls move to increase the
area of the favoured domain generating a net magnetisation in the specimen along
the direction of field H. The net magnetisation in case of antiferromagnetic materials
in presence of the magnetising field H is smaller than that of ferromagnetic material.
Susceptibility of an antiferromagnetic material is positive but small and varies with
temperature as shown in Fig. 3.20. An antiferromagnetic material behaves like a
paramagnetic material above Neel temperature T N .
The only element that shows antiferromagnetism at room temperature is
chromium (Cr) for which Neel temperature is 37 °C (310 K). Most antiferromagnetic
materials are transition metal oxides (oxides of elements whose atom has a partially
filled d-sub shell, like Cr, Mn, Fe, Co, Ni, etc.). The structure of magnetic unit cell
of MnO, a typical antiferromagnetic material is shown in Fig. 3.21a. An antiferro-
magnetic material behaves like a paramagnetic material above Neel temperature T N .
However, the variation of susceptibility of an antiferromagnetic material below Neel
temperature depends on the angle between the magnetising field H and the plane
defined by the magnetic moments of the neighbouring atoms/ions.
When the magnetising field H is parallel to the two spins S 1 and S2 of neighbouring
atoms, the susceptibility is denoted by χ P . In case H is perpendicular to S 1 S 2 , the
susceptibility is denoted by χT . Temperature dependence of χ P and χT is shown in
Fig. 3.21b.
3.7.4.1 Ferrimagnetisms
Some ceramics exhibit permanent magnetisation, called ferrimagnetism. Micro-

scopic arrangement of magnetic moments in antiferromagnetic and ferrimagnetic
materials is similar; in both the magnetic moments of adjacent atoms/ions is aligned
antiparallel to each other. In case of antiferromagnetic materials the antiparallel
magnetic moments completely cancel each other and, therefore, an antiferromagnetic
material has no magnetisation in absence of any external magnetic field. However,
in ferrimagnetic materials, the magnetic moments of the two adjacent atoms/ions
are not exactly equal: moment of one is larger than the moment of the other and,
Fig. 3.21 a Unit magnetic cell of antiferromagnetic MnO compound. b Temperature dependence
of the susceptibility of antiferromagnetic material below Neel temperature
hence, ferrimagnetic materials show a magnetisation even in absence of an external

magnetising field. Synthetic materials with ferrimagnetisms are called ferrites.
Depending on the crystal structure, ferrites may be divided into three classes:
(i) cubic (ii) hexagonal and (iii) garnet. Magnetite, Fe3 O4 is a well-known cubic
ferrimagnetic substance. The chemical composition of the material is FeO.Fe2 O3 . It
has two types of iron ions: ferrous, doubly charged Fe2+ ions and triply charged Ferric
Fe3+ ions. The compound magnetite in crystalline state has spinel structure. A unit
cell of crystal contains 56 ions out of which 24 are iron ions and 32 are oxygen ions.
In 24 Fe- ions 16 are ferric and 8 ferrous. Only iron ions have magnetic moments.
In each unit cell the Fe- ions are located in two different coordinate environments:
(a) A tetrahedral one, where the Fe-ion is surrounded by four oxygen ions and the
other (b) octahedral structure in which each Fe- ion is surrounded by six oxygen
ions. Out of 16 ferric ions 8 are in coordinate environment (a) and the remaining 8 in
environment (b). The spin or magnetic dipoles of ferric Fe- ions (Fe3+ ) in coordinate
environment (a) and (b) are directed opposite to each other and, therefore, 8-ferric
ions in both configurations cancel their magnetic moments. The resultant magnetic
moment of the molecule, therefore, arises entirely from the 8-ferrous Fe-ions which
are placed at octahedral sites. Each ferrous Fe-ion has six 3d electrons whose spin
orientations are ↑↑↑↑↑↓ Therefore, each ferrous Fe-ion carries a magnetic moment
of 4 Bohr magneton (Bm or μB ) while each ferric ion carries a magnetic moment of
5μB . Schematic representation of Fe3 O4 is shown in Fig. 3.22.
The general chemical formula for cubic ferrites is: MO.Fe2 O3 ; where M is a
divalent cation, often of Zn, Cd, Fe, Ni, Cu, Co or Mg. The general ionic formula for
such ferrite may be written as M2+ O2− (Fe3+ )2 (O2− ), where M2+ may be ions like
Mn2+ , Co2+ , Cu2+ etc. each of which may have a net magnetic moment different than
4μB . Magnetic moments in unit of Bohr magneton for some frequently used ions
are: Cu2+ (1 μB ); Ni2+ (2 μB ); Co2+ (3 μB ); Fe2+ (4 μB ); Mn2+ (5 μB ); F3+ (5 μB ).
Fig. 3.22 Spin arrangement of Ferric and ferrous ions in magnetite
Using different divalent ions and mixtures of different divalent ions it is possible to
make ferrites with desired magnetic properties. Ferrites made by mixing two or more
divalent ions are called mixed ferrite, an example is (Mn, Mg)Fe2 O4 .
Hexagonal ferrites have the general formula AB12 O19 , where A is a divalent
metal, like barium (Ba), lead (Pb) or strontium (Sr) and B is a trivalent metal like
aluminium (Al), gallium (Ga) or iron (Fe). Hexagonal ferrites have an inverse spinel
like crystal structure with hexagonal symmetry. A common example of hexagonal
ferrite is BaFe12 O19 .
A class of ferrites is called garnets that have a complicated structure which may
be represented as M3 Fe5 O12 . M in this formula stands for some rear earth ion like
yttrium (Y), gadolinium (Gd), samarium (Sm) or europium (Eu). Yttrium iron garnet
(Y3 Fe5 O12 ) is one of the frequently used garnets often denoted as YIG.
Saturation magnetisation of ferrites is not as large as that of ferromagnetic mate-
rials but their biggest advantage is that some ferrites are ceramics and excellent
electric insulators. An insulator magnetic core of ceramic ferrite in high-frequency
transformers eliminates eddy current losses.
SAQ: What are ferrites? And why they are very important? How a ferrite of desired
magnetic properties may be synthesised?
3.8 Permanent Magnetic Materials
Permanent magnets are required in all walks of life, be it big particle accelerators
used in research or tiny computer memories or fridge-magnetic stickers. Permanent
magnets are characterised in terms of the maximum energy product, i.e. the area
3.8 Permanent Magnetic Materials 183
of the largest rectangle starting from the origin that may be drawn in the second
quadrant of the B–H curve, as shown in Fig. 3.23. (BH)max tells about the magnetic
energy stored in the material per unit volume and is treated as the magnetic figure-
of-merit of the material. Maximum energy product is often measured in units of
kilo-Joule per cubic metre, (kJ m−3 ) in SI system or MGOe (Mega-gauss-Oersted)
in electromagnetic system. Further, 1 MG Oe = 7.958 kJ m−3 . Research in the field of
magnetic materials has led to almost an exponential rise in the magnitude of (BH)max
in the twentieth century starting from 20 kJ m−3 in 1900 to around 450 kJ m−3 in
2000. Increase in the (BH)max value resulted in considerable reduction in the size
of permanent magnets; for example a NdFeB magnet of 102 cc volume will contain
roughly the same magnetic energy as a brass bond lodestone of 105 cc volume; a
reduction of almost 103 in the size of the magnet.
Oldest permanent magnetic material is Lodestone, naturally occurring iron oxide
Fe2 O3 , though loadstone magnets produce low fields but they offer high resistance
to demagnetisation. Magnetic carbon steels, developed in eighteenth century, are
generally alloyed with chromium or tungsten to restrict domain wall movement and
increase coercive field. Magnets made from carbon steel have large saturation field,
order of magnitude larger than loadstone magnets, but have lower value of coercive
field. Synthetic magnets made from alloys of aluminium, cobalt and nickel, called
Alnico magnets, were first developed around 1930 and show considerably larger
values for magnetic hardness as compared to carbon steel. They also have high Curie
temperature of the order of 900 °C. Alnico is cast in a foundry. Magnets of desired
pattern may be made by using sand moulds and pouring molten magnetic material in
the mould. Alnico magnets may also be made by sintering process to form small and
accurate magnets. Alnico magnets have high operating temperature, good corrosion
resistance and long-term magnetic stability. However, their drawbacks are the high
cost on account of Cobalt and low resistance to demagnetisation. Pushing two Alnico
magnets in repulsion may demagnetise both of them.
Ferrite (Fe3 O4 ) is manufactured using powder sintering technology and exact size
tooling into range of industry standard sized disks, rings or blocks of 150 mm ×
100 mm × 25 mm size. These blocks can then be sliced into smaller magnets.
Ferrites are used extensively in loudspeakers and other security systems. The biggest
advantage of ferrite magnets is their very low cost, high resistance to corrosion and
good magnetic stability. Low level of magnetism is their only weakness.
Fig. 3.23 Maximum energy

product (BH)max defining
rectangle
A new variety of magnets, called Cobalt–Platinum magnets, developed in 1950,

their greatest advantage over the other magnetic materials is that they have excellent
resistance to corrosion and therefore they are almost exclusively used in biomedical
applications. Hard Ferrite Magnets, like BaFe12 O19 and SrFe12 O19 , are presently
the most used commercial magnetic materials for decay or two. High cost of these
magnets is their biggest weakness.
Neodymium-iron-boron compound, Nd2 Fe14 B, a magnetic material, was first fabri-
cated in 1982 by General Motors. Since then strong grade neodymium magnets
are commercially available. It is claimed that a 50 mm × 50 mm × 25 mm N52
Neodymium magnet can support vertically a steel weight of about 110 kg. The mean
value of (BH)max for Nd-magnets of different grades is around 326 kJ m−3 . These
magnets have an elaborate manufacturing process consisting of vacuum melting,
milling, pressing and sintering. These magnets have high magnetic strength, but
their weakness is their low operating temperature.
Magnetic rubber is produced by heavy loading of ferrite powder of strontium or
barium base into synthetic elastic matrix like PVC. The rubbery magnetic material
is then extruded to a desired shape or is produced in thin sheets by calendaring.
Magnetic rubber may be cut into any shape using foam cutters. The strong points of
magnetic rubber are: flexibility, ease to cut and good resistance to corrosion and the
weak points are: low magnetic flux density and low value of operating temperature.
SAQ: (BH)max , though defined in terms of the area of a specific rectangle has SI units
of kJ per cubic metre. How can one reconcile the units of area in definition
with units of volume in SI units?
Some characteristic properties of materials used in fabricating permanent magnets
are given in Table 3.4.
Some important conversion factors are: 1 Gauss = 10–4 T (Tesla); 12.54 Oe =
1kA m−1 (kilo ampere per metre; 1 MGOe = 7.958 kJ m−3 (kilo Joule per cubic
metre.
Table 3.4 Properties of some magnetic materials

Material Residual Coercive field Maximum power Maximum
magnetisation Hc (kA m−1 ) product (BH)max temperature of
Br (Tesla) T (kJ m−3 ) operation (Kelvin) K
Neodymium 13 9170 335 353
Alnico 1.25 51.00 44 773
Ceramic + 0.40 2.35 28 450
Fe2 O3
Samarium + 1.10 774.00 225 620
Cobalt
Magnetic 0.20 128.00 7 323
rubber
Solved Example SE(3.1) Calculate the value of 1 Bohr magneton (MB or μB ) in SI

units.
−19 −34
×6.626×10
qe h
Solution: By definition 1 μB = 4πm e
= 1.6×10
4×π×9.1×10−31
= 9.274 × 10−24 A m2 .
−19
Here, qe = charge on an electron = 1.6 × 10 Coulomb
h = Planck’s constant = 6.626 × 10−34 J per Hz
And m e = mass of an electron = 9.1 × 10−31 kg
Solved Example SE(3.2) A particle of mass 1 kg having charge of + 5.0 μC moving

with a constant velocity of 1.5 × 104 m/s making an angle of θ with the direction of
magnetic induction of strength 0.25 T, experiences a force of 1.7 × 10–2 N. What is
the magnitude of angle θ ?
Solution: The force F experienced by a particle of charge q moving with velocity v
in a magnetic field of induction B making an angle θ with B is given as;
F = q(v X B) = qv B sin θ
Substituting the value of q = 5.0 × 10–6 C; v = 1.5 × 104 m/s; B = 0.25 T and F
= 1.7 × 10–2 N, in the above expression, one gets
1.7 × 10−2 = 5 × 10−6 × 1.5 × 104 × 0.25 × sin θ
Or
1.7 × 10−2
sin θ = = 0.906
5 × 10−6 × 1.5 × 104 × 0.25
Therefore, θ = sin−1 0.906 = 64.96◦ .
Solved Example SE(3.3) Ferrite Fe3 O4 forms a cubic crystal with unit cell edge
length (a) of 0.840 × 10–9 m. Each unit cell of the material contains 8 Fe2+ ions with
each ion having a magnetic moment of 4 μB , and 16 Fe3+ ions that are non-magnetic.
Calculate the saturation magnetisation Bs per unit cell in units of A/m.
Solution: Total magnetic moment of a unit cell due to 8 Fe2+ ions M cell = 8 × 4 μB
= 32 μB .
( )3
Volume of the unit cell Vcell = a 3 = 0.84 × 10−9 = 0.59 × 10−27 m3 .
32 μB 32×(9.27×10−24 A m2 )
Saturation magnetisation per cell B = Mcell =
s Vcell
=
0.59×10−27 m3 0.59×10−27 m3
(on substitution of μB = 9.27 × 10−24 A m2 ).
Or Bs = 5.02 × 105 A/m.
Problems
P3.1 Two identical squares of sides 6 cm, made from conducting wire are placed
side by side on a horizontal table as shown in the figure (P3.1) and a current
of 45.0 mA is made to flow through the squares. Determine the nature and
magnitude of the force between the two squares.
ANS: Repulsive force of 78.0 × 10–10 N

P3.2 Ferrite Mn3 O4 , where Mn is the diatomic metal magnetic, forms a cubic crystal
with unit cell edge length (a) of 0.840 × 10–9 m. Each unit cell of the material
contains 8 Mn2+ ions with each ion having a magnetic moment of 5 μB , and
16 Mn3+ ions that are non-magnetic. Calculate the saturation magnetisation
Bs per unit cell.
ANS: 6.25 × 105 A/m
SA3.1 Write a note on diamagnetism giving Langvine’s theory for it. Why
diamagnetism is called the fundamental magnetism?
SA3.2 Give the major points of similarities and differences between ferromag-
netic and ferrimagnetic materials. What causes domains in these magnetic
materials?
SA3-3 What are the required magnetic properties that a material used for making
permanent magnets must have?. List some important materials used for
making permanent magnets.
SA3.4 Write the expression relating magnetic dipole moment of a charged particle
with its angular momentum in classical limits and discuss how the expres-
sion may be modified for quantum limits. Define 1 Bohr magneton and
give its value in SI units.
SA3.5 Explain why in absence of any external magnetic field H and at a temper-
ature T (a little above 0 K), diamagnetic, paramagnetic and ferromagnetic
materials all show no magnetism. Also write the order of magnitude of bulk
susceptibilities for these materials.
SA3.6 What is exchange interaction and what is its origin? Discuss the role played
by exchange interaction in case of ferro, antiferro and ferric magnetic
materials.
SA3.7 What are ferrites? Explain how magnetic material with desired value of
saturation magnetisation may be fabricated using mixtures of ferrites.
SA3.8 Draw a typical B–H curve for a ferromagnetic material indicating important
characteristics of the curve. What does the area of the B–H curve represents?
Define (BH)max and discuss the significance of this parameter.
SA3.9 Explain how the application of an external magnetic field H in case of metals
that have free electrons, causes an imbalance in the number of electrons
with opposite spins that leads to a lower value of magnetisation in the
direction of field H.
SA3.10 Explain what is magnetic anisotropy and what causes it.
SA3.11 What are the distinguishing features of magnetically hard and soft mate-
rials? Briefly outline applications of the two types of these materials. What
are ceramic soft magnetic materials and where are they used.
SA3.12 What is meant by singlet and triplet states of spin alignment? Discuss their
role in magnetism.
SA3.13 Draw rough sketches for the variation of bulk susceptibility χ and (1/χ )
with temperature for different types of magnetic materials and define the
Curie and the Neel temperatures.
SA3.14 Write a note on ferrimagnetism giving some examples.

Note: Some of the multiple choice questions may have more than one correct alter-
native. All correct alternatives must be marked for the complete answer of such
questions.
MC3.1 Biot–Savart law in mathematical form may be given as;
{ −
→ { −
→ { −
→
(a) B→ = μ4π0 I dlrX2 r̂ (b) B→ = μ4π0 I dl|r|×→
3
r
(c) B→ = 2πI dl X r̂
r2
(d) B→ =
I
{ −
→
dl ×→r
4μ0 π r2
ANS: (a), (b)

MC3.2 Magnetic dipole moment of an object is equal to the
(a) Minimum torque experienced by the object in an external magnetic
field
(b) Minimum torque experienced by the object in an external magnetic
field of strength unity
(c) Maximum torque experienced by the object in an external magnetic
field
(d) Maximum torque experienced by the object in an external magnetic
field of strength unity
ANS: (d)
MC3.3 A circular loop of wire carries a current I in clock-wise direction. The

direction of magnetic field at a point P inside the loop is
(a) Left to point P (b) right to point P (c) points out of the page (d) points
into the page
ANS: (d)
MC3.4 1 Bohr magneton is
(a) 9.274 × 10−24 J/T (b) 9.274 × 10−21 erg/gauss (c) 9.274 × 10−21 eV/T
(d) 5.7883 × 10−5 eV/T
ANS: (a), (b) and (d)
MC3.5 M L 3 T −1 Q −1 ; and N M T −1 are dimensional formulas for
(a) Magnetic moment (b) Magnetisation (c) Bohr magneton (d) Magnetic
induction
ANS: (c)
MC3.6 ‘g’-factor for the spin motion of electron is around
(a) 1 (b) 2 (c) 3 (d) 4
ANS: (b)
MC3.7 Stick out the incorrect alternative from the followings:
Bulk susceptibility of diamagnetic materials depends on
(a) the number of electrons per atom (Z),
(b) number density (N) of atoms (atoms per unit volume)
(c) temperature of the specimen T
(d) Intensity of the magnetising field H
ANS: (c), (d)
MC3.8 Langevin’s theory for paramagnetism is not applicable to.
(a) metal ions (b) hydrated salts of transition metals (c) dilute gases (d)
diamagnetic substances
ANS: (a), (d)
MC3.9 Exchange interaction in ferromagnetic materials is the result of
(a) atomic dipole–dipole interaction and Pauli’s exclusion principle
(b) lattice-dipole interaction and Pauli’s exclusion principle
(c) spin-dipole interaction and Pauli’s exclusion principle
(d) interaction between electron clouds and Pauli’s exclusion principle
ANS: (d)
MC3.10 Best material for the core of a transformer is

(a) diamagnetic (b) paramagnetic (c) antiferromagnetic (d) ceramic ferrites
ANS: (d)
MC3.11 Which of the following are intrinsic magnetic properties of a crystalline
magnetic solid?
(a) Curie temperature (b) saturation magnetisation (c) crystalline magneto
anisotropy (d) bulk susceptibility χ
ANS: (a), (b), (c)
LA3.1 What are the characteristics of ferromagnetic materials? Which interactions

produce domains in magnetic materials? Discuss in details Weiss theory of
ferromagnetism.
LA3.2 Define magnetic permeability and susceptibility. What are paramagnetic
materials? Discuss with necessary detail Langevin’s theory for paramag-
netism and hence bring out the importance of Curie temperature.
LA3.3 Explain what causes diamagnetism in materials the atoms/ions of which
have no magnetic moment. Give a detailed account of Langevin’s theory for
diamagnetism.
LA3.4 What are the distinguishing features of diamagnetic, paramagnetic, ferro-
magnetic, antiferromagnetic and ferrimagnetic materials?. What forces/
interactions are responsible for these differences? Discuss the origin of
exchange interactions and the role they play in magnetism.
Chapter 4
X-rays, Dual Nature of Matter, Failure
of Classical Physics and Success
of Quantum Approach
Objective
After reading this chapter the reader is expected to: (i) understand the method of
producing X-rays, their properties and applications, (ii) be able to grasp the concepts
of the dual nature of matter, energy and matter waves, (iii) appreciate the pitfalls of
classical theories of physics in explaining some important experimental observations
and how the quantum approach is able to explain them.
4.1 Introduction
The foundation of classical physics, considered to be propounded in pre-1900 era,

rests with the laws of motion given by Newton, theory of electromagnetism given
by Maxwell and the laws of thermodynamics. During this period, the X-rays were
discovered and were found to have many properties that initially indicated their quite
abnormal behaviour. The X-ray photons were found to be like light photons except
they contain so much energy, that they penetrate thin opaque sheets. Though the X-
rays have sufficient energy to penetrate thin films of material, but objects of higher
densities were found to block X-rays, casting shadows on screen placed in front of
these. It was also realised that X-rays, like the visible light, are electromagnetic in
nature. During the same period a major and revolutionary idea of matter waves was
proposed by de Broglie. The idea that moving material particles have an associate
wave, i.e. the dual nature of matter, was experimentally confirmed by Davisson and
Germer. However, there was some phenomenon, for example, the energy distribution
in blackbody radiation spectrum, interaction of light with matter, the heat capacity
of gases, stability of atom, its line spectra, etc., that could not be explained by
the available classical physics. In this chapter, it is proposed to briefly discuss the
production and properties of X-rays and the experiments to establish the dual nature
of matter. Failure of classical physics in explaining some important phenomenon is
https://doi.org/10.1007/978-3-031-32084-2_4
192 4 X-rays, Dual Nature of Matter, Failure of Classical Physics and Success …
discussed in the later part of the chapter. Quantum mechanical explanations of these
anomalies are also included, towards the end of the chapter.
4.2 Discovery, Production and Properties of X-rays
The discovery of X-rays, that was accidental, was announced by Wilhelm Conrad
Roentgen, in December 1895. W. C. Roentgen, Physics professor at Wurzburg
University, Bavaria, Germany, was studying the properties of cathode rays that are
emitted when an electric discharge is made to pass through the two electrodes of a
cathode tube filled with some gas at low pressure. Roentgen was specifically inter-
ested in whether cathode rays could pass through glass body of the cathode-ray
tube. His cathode-ray tube was covered by thick black paper on all sides, but he was
surprised to see that an incandescent green light nevertheless escaped and projected
onto a fluorescent screen placed nearby. Roentgen found that these mysterious rays
could penetrate through most substances and cast their shadows on the screen. Since
the exact nature and properties of these rays were not known at that time Roentgen
called them X-rays, the term ‘X’ usually used to describe the unknown quantity.
Roentgen also found that the X-ray was capable of passing through human’s tissues
leaving the shadow of bones. Almost immediately X-ray’s uses as a diagnostic tool
to detect bone fractures become a common medical practice.
4.2.1 Production of X-rays
The X-rays are generally produced in a specially designed vacuum tube invented by
William Crookes which is often called the discharge tube or a cathode-ray tube. A
typical sketch of an X-ray tube is given in Fig. 4.1.
As shown in the figure, the X-ray tube consists of an evacuated glass tube which
is fitted with two electrodes. The electrode which is kept at a negative potential is
called cathode and the other electrode kept at a higher positive potential, the anode.
A potential difference of the order of few tens of kV is maintained between the
two electrodes. Often a heater element is also attached to the cathode which may
be connected to a low-voltage source. When a current is passed through the heating
element, the temperature of cathode increases and thermionic emission of electrons
from cathode takes place. The emitted thermo-electrons get repelled by the nega-
tive potential at cathode and are attracted by the positive anode. Thus electrons get
accelerated and impinge with high speed on to the anode when they get decelerated.
Since accelerated/decelerated charged particles (electrons in this case) emit electro-
magnetic radiations, decelerated electrons in the X-ray tube emit electromagnetic
radiations in the form of X-rays. The X-rays generated due to deceleration of elec-
trons are termed as continuous, soft or Bremsstrahlung X-rays. Bremsstrahlung
X-rays consist of X-ray radiations of all energies starting from a minimum energy
4.2 Discovery, Production and Properties of X-rays 193
Fig. 4.1 Schematic diagram of an X-ray tube. X-rays are produced when energetic electrons
impinge on the anode material
E min to a maximum energy E max , hence the name continuous X-rays. X-rays are also
produced when high-energy electrons hit the atoms of the anode material and shift the
atomic electrons to their higher energy states, thus exciting the anode atoms. Since
atoms cannot remain in excited states for long, excited atoms of the anode revert
back to their ground states emitting X-rays. X-rays produced by the de-excitation
of excited atoms are called characteristic X-rays and their energy depends on the
atoms of anode material. Characteristic X-rays have X-rays of some very definite
wavelengths that depend on the nature of the target atom and are found superimposed
as sharp lines on the continuous background of Bremsstrahlung X-rays.
Since their inception, the X-ray tubes have evolved and undergone several
changes/improvements. A modern cold cathode X-ray tube is shown in Fig. 4.2.
As shown in the figure, the modern X-ray tube contains an anticathode or target
opposite the concave-shaped cathode, further the cathode is not heated and electrons
are not emitted by cathode. The X-ray tube is filled by some inert gas like argon at
low pressure. A spark plug is used to ignite the inert gas which on ionisation produces
electrons. Electrons produced by the inert gas are accelerated between the concave
cathode and the target or the anticathode. Any desired material may be attached to
the anticathode in order to produce characteristic X-rays of that material. Concave
cathode focuses the electron beam to a point on the target which helps in generating
focused X-rays that produce sharp images of dens materials like bones, etc. These
X-ray tubes do not have any heater at its cathode; therefore, they are referred as cold
cathode tubes.
Soon after the discovery of X-rays, efforts were made to study the properties
of these rays. Experimental observations revel that X-rays have a high penetrating
power, travel in straight line and cannot be deflected by electric field or the magnetic
field. When the X-rays are incident on the photographic plates, they are found to
Fig. 4.2 Modern X-ray tube
blacken the film. When X-rays pass through a gaseous medium, they ionise it and
can also cause photoelectric emission similar to when light is incident on metallic
cathode. The X-rays have broadly been categorised into two categories as continuous
X-rays and discrete/characteristic X-rays.
4.2.2 Continuous X-rays
Continuous, Bremsstrahlung, white or soft X-rays are produced when energetic

electrons, accelerated by the potential difference between the anode and the cathode
of an X-ray tube, get retarded on hitting the target (anode or anticathode) which is
generally of some metal. Electrons approaching a metal target interact with atoms
of the target in two ways: (i) they are repelled by the negative electron cloud of
target atom and (ii) they get inelastically scattered by the positively charged nucleus
of the target atom. In both these interactions, incident electrons lose energy. The
amount of energy lost by each incident electron depends on several factors including
its direction of motion, distances from the electron cloud and nuclei of target atoms
and the angle of scattering. Different electrons of the incident electron beam lose
different amounts of energies on interacting with target atoms. The energy lost by
each electron is converted into an X-ray. In this way, a bunch of X-rays having a
random energy distribution (or wave length) is produced when a beam of accelerated
electrons hit the target electrode in an X-ray tube.
Bremsstrahlung is a word of German language, which may be broken into two
German words: Brems, which means break (stopping), and strahlung, that means
Fig. 4.3 Deceleration of incident electron by the coulomb field of the electron cloud of the target
atom produce Bremsstrahlung X-ray
radiations. Bremsstrahlung, therefore, stands for the radiations that are emitted by
the stopping of electrons.
The process of X-ray emission as a result of deceleration of electrons is called
Bremsstrahlung. Figure 4.3 shows how an incident electron gets retarded in the
Coulomb field of the electron cloud of the target atom, and the difference of kinetic
energy ΔE at points A and B is converted into an X-ray.
A figure depicting the emission of X-ray as a result of change of velocity of
electron due to scattering by the positively charged nucleus of the target atom is
shown in Fig. 4.4.
Fig. 4.4 Deceleration and deflection of an energetic electron by a positively charged nucleus,
resulting into emission of X-ray photon
X-rays, like other components of electromagnetic radiation spectrum, exhibit both

the wave and particle nature. The X-ray energy quanta are called photon. X-rays are
characterised either by energy E x , frequency ν or the wavelength λ, which are related
to each other through following expressions:
hc c E hν
E x = hν = ; ν= ; p= = (4.1)
λ λ c c
Symbols h, c and p in above expressions stand respectively for Planck’s constant,
speed of light and the linear momentum of X-ray photon.
Intensity distributions of continuous X-rays as the function of wavelength and as
the function of photon energy E x or frequency γ are shown in Fig. 4.5. It may be
noted in Fig. 4.5a that X-ray distribution curve has a cut-off wavelength, denoted
by λmin while on the higher wavelength side the curve extends almost up to infinity.
Minimum wavelength λmin corresponds to the X-ray of the highest energy, similarly
X-rays of longer wavelengths have smaller energies extending to almost zero energy.
In Fig. 4.5b, the X-ray of highest energy has the largest frequency γ max and on the
lower frequency side the distribution curve extends almost to the origin. An important
property of continuous X-rays is that the cut-off point (λmin or γ max ) depends only
on the potential difference V between the anode and the cathode of the X-ray tube
and it does not depend on the target material.
The minimum wavelength λmin and maximum energy or maximum frequency
νmax of emitted X-ray corresponds to an incident electron losing all of its energy in
a single collision and radiating it away in the form of a single X-ray photon. If we
assume that total kinetic energy (K.E.) of the electron is converted into energy (hν)
of the X-ray photon, then
Fig. 4.5 Intensity distribution of continuous X-rays as a function of a X-ray wavelength, b X-ray
energy or frequency
Kinetic energy lost by the electron ΔE = Energy hν of emitted X-ray
Or

c
ΔE = hνmax = h (4.2)
λmin
hc
hνmax = (4.3)
λmin
The λmin may be called as the cut-off wavelength, which will mainly depend on
the value of accelerating voltage V applied across the anode and cathode. Thus,
hc
hνmax = = eV (4.4)
λmin
And
hc
λmin = (4.5)
eV
When one substitutes the values of h, c and e in Expression (4.5), one gets the
simplified relation between λmin (in units of Å = 10−10 m) and the voltage V between
the anode and cathode of the X-ray tube as,
12.398 × 103
λmin in units of Å = (4.6)
V (volts)
The total X-ray energy emitted per second (intensity) I depends on the atomic
number Z of the target atom, the electric current i passing through the two electrodes
of the X-ray tube and the potential difference V between the electrodes. Intensity I
may be written as:
X-ray intensity I = Ai Z V m (4.7)
Here A is a constant and m is also a constant that has the value of ≈ 2 for most
of the metals.
The dumbbell-shaped spectrum of continuous X-rays has a broad peak which
corresponds to the wavelength of the X-rays of maximum emission, the X-rays
which have the highest number in intensity distribution. With the increase of the
voltage V between the electrodes of the X-ray tube, the peak of maximum emission
and the cut-off wavelength λmin both shifts towards the shorter wavelength and the
total area of the intensity curve also increases as shown in Fig. 4.6. This happens
because of the increase in the kinetic energy of electrons with voltage V, resulting
in larger number of electron interactions in which more energy is lost in the form of
X-rays, and because current of the X-ray tube i also increases with V.
Fig. 4.6 Variation of the continuous X-ray spectrum with the voltage between electrodes
4.2.3 Characteristic X-rays
Continuous component of X-rays is produced from the deceleration of high-energy

electrons in the X-ray tube by Coulomb fields of negative electron cloud and the
positively charged nucleus of the target atom. The electrons in the X-ray tube are
generally accelerated by high voltages of the order of few tens of kilovolts, have
sufficient energy to initiate two other processes. (i) They may excite the atoms of
the target material by shifting an electron from the lower energy state to a state of
higher energy. (ii) High-energy electrons may ionise the atoms of the target material
by pushing out some electrons from the atom. Both the excited and the ionised target
atoms are unstable, they cannot remain in excited or ionised state for long and in most
cases the unstable target atoms revert back to their ground states within 10−9 –10−7 s.
Characteristic X-rays are emitted by the atoms of the target material when they
de-excite from the ionised or excited states. The wavelength and intensity of these
X-rays depend on the type of the atom of the target material, and therefore, they are
called characteristic X-rays. The energy (or wavelength/frequency) and intensity of
characteristic X-rays change with change of the target material attached with anode
or anticathode of the X-ray tube (Fig. 4.7).
Characteristic X-rays are produced when an incident electron of high energy
interacts with a bound electron (say K-shell electron) and eject the K-shell electron
out of the atom, creating a vacancy of electron in K-shell. Obviously, this is possible
only if the energy of the incident electron is more than the binding energy of electron
in K-shell. This is ionisation of the target atom. The ionised atom is unstable and
Fig. 4.7 Representation of X-ray emission during the transition of an electron from L- to K-shell
may revert back to a state of lower excitation by the transfer of an electron from the
next higher shell, the L-shell. The difference in the energy of the ionised atom with
electron vacancy in K-shell and when electron vacancy is in L-shell, is emitted in the
form of an X-ray. This X-ray may be denoted as X L→K and is called the K α X-ray.
Similarly, the vacancy in the K-shell may be filled by an electron from the M-shell
(instead of from L-shell) producing a X M→L -ray which is designated as K β X-ray.
In this way X-rays of different energies may be produced by the de-excitation of an
ionised target atom. In case the high-energy electron in the X-ray tube creates an
electron vacancy in L-shell of the target atom (instead of the K-shell), then in that
case electrons from M, N and other higher shells may fill the electron vacancy in
L-shell producing L α , L β , L γ . . . X-rays. Figure 4.8 shows the electron energy level
diagram of any general target atom and the transitions corresponding to various X-
rays. In the energy level diagram it may be noted that the energy difference between
successive energy levels decreases as one goes higher in energy. The maximum
energy difference occurs between the K- and L-shells, and therefore, the energy of
the K α X-ray is equal to (E L − E K ), where E L and E K are respectively the energy of
the L-shell and of the K-shell. In an actual case, there are large number of high-energy
electrons in the X-ray tube that hit the target almost simultaneously and ionise many
target atoms producing vacancies in different shells of different atoms, as a result
many characteristic X-rays like K α , K β , L α …, etc. are emitted by target atoms.
Since atoms of different materials have different energy level diagrams, the ener-
gies of characteristic X-rays are different for different materials. The characteristic
X-ray spectrum may therefore be treated as a finger print of the atom and is often
used to identify different atoms. The characteristic X-ray spectrum is always found
to be superimposed as sharp peaks on the background of continuous X-ray spectrum
as shown in Fig. 4.9. Since the width of characteristic X-ray peaks is very small as
Fig. 4.8 Schematic energy

level diagram of a target
atom and transitions
corresponding to different
characteristic X-rays
compared to the broad peak of continuous spectra, characteristic peaks are called
characteristic lines.
Most X-ray tubes operate at voltages of the order of 50 kV or less. Therefore,
the maximum energy of continuous X-rays may be at the most 50 keV or a little
less, when an electron accelerated to 50 keV energy loses all its energy in a single
event producing an X-ray of 50 keV. As such the continuous X-ray spectra from
such an X-ray tube will terminate at 50 keV energy. The maximum energy X-ray
in characteristic spectra will correspond to K α line of the target atom. For medium-
weight target elements the energy of K α line lies in the range of 50–80 keV energy.
Evidently, in such cases, the characteristic X-ray lines are found to be superimposed
on the high-energy tail of the continuous spectra, as shown in Fig. 4.9.
As already mentioned, the characteristics (wavelength of maximum emission,
intensity and cut-off or minimum wavelength) of the continuous X-ray spectrum
depends essentially on the magnitude of the voltage between the anode and the
cathode of the X-ray tube; however, the wavelengths or frequencies of characteristic
Fig. 4.9 Peaks of

characteristic X-ray
spectrum superimposed on
the background of
continuous X-ray spectrum
X-ray lines depend only on the atoms of the target material and do not change with
the voltage between the two electrodes of the X-ray tube.
SAQ: What is the basic origin of X-rays and in what respect their origin is different
from that of gamma rays?
SAQ: Which characteristic X-ray of a given atom will have highest energy and
why?
4.2.4 Mosley’s Law
With the discovery of large number of elements, attempts were made to arrange
elements in some order. First such attempt was made by John Dalton in 1803 when he
arranged elements according to the increasing atomic weight. Later, it was observed
that groups of elements exhibit similar chemical properties suggesting the presence of
recurring patterns of chemical behaviour. Around 1870 Dimitri Mendeleev developed
what is called the periodic table of elements, where elements were largely placed
according to their atomic weights and numbered consecutively. In this periodic table
no physical meaning or significance was attached to the sequence number of the
element. However, some anomalies were found in this arrangement of elements in
the table; for example, the atomic weight of Cobalt was higher (58.93) than that of
Nickel (58.69), but its chemical properties suggested that it should be placed before
Nickel in the periodic table. Similarly, the atomic weight of argon was larger than
that of potassium, and the atomic weight of tellurium was greater than that of iodine,
but the chemical properties of both these elements suggested that they in spite of
their heavier weights should precede the corresponding lower weight partner. These
anomalies indicated that atomic weight is not the correct criteria for numbering
elements in the periodic table. Soon it was realised that numbering of elements in
periodic table should be done according to the number of electrons in the atom of the
element which is equal to the number of units of positive charge on the nucleus of the
atom and is denoted by Z, the atomic number. A final and definitive resolution of this
anomaly was achieved by H. G. Moseley, an English physicist, who in 1913 published
a research paper based on his analysis of characteristic X-ray spectra from several
elements showing that the frequencies of characteristic X-ray lines are proportional
to the squares of whole numbers that are equal to the atomic number plus a constant.
Moseley used Bohr’s atomic model for the analysis of experimentally observed
characteristic X-ray spectra from many elements. In order to appreciate Moseley’s
analysis it is required to re-visit Bohr’s model of the atom.
In Bohr’s model of the atom, it is assumed that the electron moves in a circular
orbit around the positively charged point nucleus, balancing the centrifugal force
by the attractive Coulomb force between the oppositely charged electron and the
nucleus. The breakthrough in Bohr’s model was the quantization of electrons’ angular
momentum. Using this model Bohr derived the following formula for the wave-
length of electromagnetic radiations emitted during transitions of electron between
quantized electron energy states in hydrogen-like atom.

1 1 1
= R 2 − 2 Z2 (4.8)
λ nf ni
Here, n i , n f respectively, denote the angular momentum of the electron in the

initial and final states (n = 1 for K-shell, 2 for L-shell, 3 for M-shell and so on), Z
the units of positive charge on the nucleus, i.e. the atomic number of the element. R
is the Rydberg constant. Two different values for this constant may be used (a) if it is
assumed that the nucleus is of infinite mass as compared to the mass of the electron,
the value R∞ is used. This assumption is true for heavy elements (b) in case it is
assumed that the mass of nucleus is finite of value M, then the value given by R is
used which is more appropriate for light elements.
m e · e4
R∞ = (4.9)
4π cℏ3
and

Mm e ·e4
R= (4.10)
M + m e 4π cℏ3
Symbols used in above expressions have following meanings: M = mass of the

nucleus; m e = mass of electron; e = unit of charge = charge of electron; c = velocity
of light and ℏ = rationalised Planck’s constant.
Though Bohr’s model was developed only for the hydrogen atom with Z = 1,
but it worked rather well in case of singly ionised helium atom for which Z = 2. It
was stipulated that if the model is applicable for the cases of Z = 1 and Z = 2, then
it may also be valid for heavier elements with Z > 2. And if so, it will be possible
to order transitions between electron shells as a function of Z alone, providing an
unambiguous meaning to the atomic number Z and ordering of elements according
to Z in periodic table.
In case of K α transition the frequency να of the transition between nf = 1 and
n = 2 shells may be obtained from Eq. (4.8) of Bohr model as:

c 1 1
νk = = Rc 2 − 2 Z 2
λ 1 2
Or
3Rc 2
νk = Z
4
And
/
√ 3Rc
νk = Z (4.11)
4
With this background, we now discuss Moseley’s work and his observations.
A systematic examination of the characteristic X-ray radiations was carried out by
Moseley for large number of elements from Aluminium to gold. He recorded the X-
ray spectra for these elements on a photographic
√ plate. The analysis of the spectra was
done, and the square root of the frequency ( ν) of the particular characteristic X-ray
radiation versus the ordinal of the element’s position in the periodic table, which we
for the present denote by number N, was plotted. Figure 4.10 shows a representative
graph where square root of the frequency of K α lines of different elements is plotted
against the serial number N of the element in periodic table. It was observed that
with the increase in the position of the element in the periodic table, the root of the
frequency of the emitted radiation increases monotonically and the graph may be
fitted with a straight line. Similar graphs for K β lines and for other lines of K and
the L-series were plotted, and it was observed that for each characteristic X-ray the
data points fall on a straight line. These straight lines may be characterised in terms
of their slope which may be denoted by ‘a’ and their intercept ‘b’ on the X-axis.
Moseley observed that the intercept ‘b’ for all members of a series is same, while
each straight line has a different value of slope ‘a’ as shown in Fig. 4.10. √
On the basis of his findings Moseley reached to the conclusion that the ν for
different series may be written as:
Fig. 4.10 Plot between the square root of the frequencies of K and L lines and the serial number
N of the element in the periodic table
√
νkα = akα (N − bk )
√
νkβ = akβ (N − bk )
(4.12)
.....................
.....................
√
ν Lα = a Lα (N − b L )
√
ν Lβ = a Lβ (N − b L ) (4.13)
.....................
√
ν Mα = a Mα (N − b M ) (4.14)
That in general may be written as:

√
ν = a(N − b) (4.15)
The constants a and b in above equations are respectively called the proportion-
ality constant and the screening constant. Moseley also obtained the values of the
proportionality and screening constants for different series from his experimental
plots and found that they may be given in terms of the Rydberg constant R, as given
below,
/ /
3Rc 8Rc
akα = and bk = 1; akβ = and bk = 1,
4 9
/ /
5Rc 3Rc
a Lα = and bk = 7.4; a Lβ = and bk = 7.4.
36 16
Substituting the above values, the root of frequency of K α -series transitions may
be written as:
/
√ 3Rc
νk = (N − 1) (4.16)
4
Expression (4.16) obtained by Moseley from analysis of K α -lines in characteristic

X-ray spectrum of many elements is an empirical relation that may be compared
with the theoretical expression (4.11) based on Bohr model for hydrogen-like atoms.
Except of the screening constant (bk = 1) in Moseley’s empirical relation, the two
expressions are identical if one takes N = Z. Moseley argued that the screening
constant originates from the fact that the effective charge Z of the nucleus is reduced
by the factor b due to the screening or shadowing of the nuclear charge by the electrons
left in a given orbital after ionisation. For example, the K-shell of an atom may have
at the most two electrons, and if one electron is removed for creating a vacancy in
Kth shell, the remaining one electron in the K-shell will screen the nucleus and the
effective charge of the nucleus will appear to be (Z − 1) units.
Moseley assumed that there must be some physical attribute of atoms of the
periodic table that increases in a regular fashion by some fixed amount, from one
element to the next one. He postulated that this can be the charge (Ze) of the nucleus of
the atom which is screened by the negatively charged electrons remaining in orbitals
after the creation of vacancy in a shell.
According to Mosley, the ordinal or serial number N of the element’s position
in the periodic table is equal to the number Z , related to the positive charge (Z e)
carried by the nuclei of the element. The number Z is referred to as the atomic
number of the element and is exactly equal to the number of protons in the nucleus
of the atoms of the materials emitting X-rays. It may be mentioned here that before the
investigations carried out by Mosley, the arrangement of elements in the periodic table
was in the ascending order of their atomic weights and on the basis of their chemical
properties. The Mosley’s results could provide a direct method of determining the
atomic number of the elements and helped in removing the discrepancies in the
periodic table arrangements. As already mentioned, initially the positions of the
transition metals Cobalt (Z = 27) and Nickel (Z = 28) were determined on the basis
of the ascending order of their atomic weights as Ni = 58.71 and Co = 58.93 were
changed. In the same way some empty positions for the still undiscovered elements
were filled. For example, new elements Hafnium (Z = 72), Technetium (Z = 43)
and Rhenium (Z = 75) were discovered as there were missing gaps at these values
of atomic numbers.
It may be noted that the difference in the magnitudes of proportionally constant
‘a’ for different members of a given series, for example, between aKα , aKβ , aKU ,
etc., is very small, and therefore, curves for different members of a series shown in
Fig. 4.10 often merge together, if the resolution of the graph is not good.
It might be of interest to know that Henry Gwyn Jeffreys Moseley was born in
Weymouth, Dorset, England, on November 23, 1887. After is education at Trinity
College, Oxford, he was appointed lecturer in Physics at Rutherford Laboratory,
University of Manchester in 1910. His initial research work was on radioactivity, but
later he carried out detailed studies on X-ray spectra and in 1913–14 published his
famous Moseley law, which paved way to uniquely determine the atomic number of
elements and their positioning in the periodic table. In 1914 he was drafted in army
during the First World War and was shot in the head by a Turkish sniper at the battle
of Suvla Bay. He died at the young age of 27 years.
SAQ: What is the physical significance of Mosley’s law?
SAQ: Calculate the magnitude of proportionality constant aMβ .
4.2.5 X-ray Diffraction
X-rays, like other components of electromagnetic radiations, undergo diffraction.

Diffraction occurs when a wave disturbance encounters an obstacle, object or aper-
ture and bends or spreads round the edges of the obstacle. In case the incident
radiation is monoenergetic of single wavelength, secondary waves originating from
different points of the incident wave front interfere with each other giving rise to
an interference pattern. The interference pattern produced by the diffraction of elec-
tromagnetic waves is characteristic of the obstacle, that is the interference pattern
may be treated as the finger print of the diffracting obstacle. The resolution of the
diffraction–interference pattern depends on the wavelength of the electromagnetic
radiations; radiations of shorter wavelength have better resolution and may resolve
structural details of objects/obstacles comparable in dimensions with the wavelength
of the radiation. X-rays are electromagnetic radiations with wavelengths in the range
of nanometer (10−9 m) and are therefore capable of deciphering structural details
of crystals where atoms or group of atoms are arranged at regular and repetitive
distance in nanometer dimensions. X-ray diffraction (XRD) is now a well-establish
and routine procedure to determine the lattice parameters, arrangement of individual
atoms in a single crystal and the phase analysis in case of polycrystalline materials
and compounds.
It may be recalled that crystals are made up of unit cells which are the simplest
repeating structure of a crystalline solid. Arrangement of unit cells results in a crystal
lattice, which is a specific three-dimensional arrangement of unit cells. Incident X-
rays on diffraction from atoms of the crystal lattice produce a distinct interference
pattern characteristic of the atomic arrangement in the crystal lattice (Fig. 4.11).
Fig. 4.11 Two-dimensional lattice with groups of crystal planes

In order to understand X-ray diffraction by crystalline materials, it is required

to understand crystal planes. Crystal planes are imaginary planes inside a crystal
which have large number of atoms. In the same crystal lattice there may be several
groups of crystal planes with different orientations. The important property of a
group of crystal planes is that all crystal plains in a group are parallel to each other
and are separated from each other by a fixed distance generally denoted by ‘d’. The
separation between successive members of a group of crystal plane, the orientations
of possible groups of planes, density of atoms in a group of particular orientation, etc.
are all depend on the crystal structure. As you might be knowing that there are seven
basic crystal systems (Triclinic, Monoclinic, Orthorhombic, Trigonal, Hexagonal,
Tetragonal and Cubic systems) and that for each system there are specific values of
crystal plane parameters. Therefore, it is possible to determine various parameters
of crystal planes using X-ray diffraction and to identify the crystal system. Since the
density of atoms in a crystal plane is large, it is reasonable to assume that each point
on this plane acts as if an atom is placed there. This assumption is rigorously valid
when the crystal plane is scanned by X-rays of wavelength of the order of nanometer
(10−9 m), while inter-atomic distance in crystal planes is much smaller of the order
of 10−10 m. With this assumption it is possible to replace atoms in the plane by
the plane itself and consider the interaction of incident X-rays by the crystal plane
instead of with individual atoms.
Let us now consider a parallel beam of monoenergetic X-rays incident on a crystal.
Individual X-ray on interaction with individual atom may be scattered or diffracted in
any direction; similarly, X-rays falling on a group of crystal planes are diffracted in all
possible directions. However, X-rays diffracted by two successive crystal planes of a
group after diffraction in a particular direction may undergo constructive or destruc-
tive interference creating a patterns that is specific to the crystal lattice. Conditions for
constructive interference between X-rays diffracted (or scattered) by two successive
crystal planes of a group are derived in the following.
Figure 4.12 shows two monoenergetic (wavelength λ) X-rays, denoted as 1 and
2, parallel to each other impinging on a family of crystal planes making an angle
θ with planes. Let us consider two consecutive planes, plane-1 and plane-2 which
scatter (or diffract) these rays as ray-3 and ray-4, respectively. The scattered rays 3
and 4 will undergo constructive interference if the phase difference between the rays
(1 + 3) and rays (2 + 4) is either zero or is an integer multiple of 2π, i.e. if the phase
difference between rays (1 + 3) and rays (2 + 4) = n(2π ), where n is an integer
including 0, the rays will give a maximum of intensity. Since a phase difference of
2π is equivalent to a path difference of λ, where λ is the wavelength of X-rays, the
condition for constructive interference in rays-3 and 4, becomes:
Path difference between rays (1 + 3) and rays (2 + 4) = nλ

(4.17)
where n = 1, 2, 3, . . .
To calculate the path difference between rays (1 + 3) and rays (2 + 4), we drop
perpendiculars AC and AD from point A on ray-2 and ray-4, respectively. As is clear
Fig. 4.12 Interference of monoenergetic X-rays scattered by successive crystal planes
from the figure, ray-2 travels a distance CB more than ray-1 and ray-4 travels the
distance BD more than ray-3. Therefore, the total path difference (Δ-path) between
the incident and scattered rays is:
Δ-path = C B + B D = d sin θ + d sin θ = 2d sin θ (4.18)
Hence, constructive interference between X-rays scattered from consecutive

crystal planes will take place when:
2d sin θ = nλ where n = 1, 2, 3, . . . (4.19)
The father and son W. H. Bragg and Lawrence Bragg for the first time derived
the above condition of constructive interference between electromagnetic radiations
scattered by consecutive members of a family of crystal planes in 1913; the condition
is called Bragg’s law.
It is easy to verify in Fig. 4.12 that the diffracted rays (2 and 4) are rotated from
their original direction (incident rays 1 and 2) by angle 2θ.
SAQ: The phenomenon of diffraction of X-rays indicates their particle nature or
wave nature?
4.2.6 Some Application of X-rays
X-rays are extensively used in medical world for detecting fractures in bones,
detecting breaks/tearing of ligaments, sterilising of medical instruments, cloths,
bandages, etc. High-energy X-rays are now used for exploring underground structures
and are providing valuable information of archaeological interest without excavation

and digging. Another very important field of application of X-rays is the character-
isation of crystalline solids. X-ray diffraction patterns may give information about
the size of unit cell, phases involved in the structure and many other characteristics.
Powder XRD, the X-ray diffraction pattern from powdered crystals has become a
powerful tool for crystal structure studies. Some details of this technique are provided
here.
In powder X-ray diffraction, a diffraction pattern is obtained from the powder
of the material, rather than an individual crystal. Powder X-ray diffraction is easier
and simpler than the crystal diffraction as no single individual crystal is required.
PXRD is characterised by high sensitivity, reliability, depth profiling, easy sample
preparation, convenient procedure, fast speed and effective resolution. Further, the
data obtained from PXRD may be used both for qualitative as well as for quantitative
analyses. However, there are some disadvantages in this method. Firstly, one uses
X-rays that are harmful for human being, and there is always some chance of leakage
and unwanted exposure of the person carrying out measurements. Secondly, and more
importantly, analysis of data requires standard references to match the experimental
pattern. Some of these difficulties have now been taken care by turn-key PXRD
systems commercially available with automatic data handling, comparing with in-
built standard references and finally providing the required information in the output.
In commercially available PXRD systems, samples are generally in the form of finely
divided powders and the pattern is generated by the diffraction of monoenergetic
X-ray from surfaces coplanar to the sample holder.
Figure 4.13 shows the simplified sketch of a power XRD system. In this configu-
ration, the X-ray source and the detector are both rotated by the same angle to look
for the intensity of diffracted X-rays. X-ray detector in old version of the system use
to be a photographic plate; however, in modern XRD systems scintillation spectrom-
eters are used to record the intensity of diffracted X-rays. Scintillation detectors are
much more efficient, have better energy resolution which helps in deciphering the
crystal structure. A typical spectra of diffracted X-rays as a function of the angle θ
are shown in Fig. 4.14.
Fig. 4.13 Simplified sketch of a possible configuration of X-ray source, powder sample holder and
X-ray detector for PXRD scanning. In this configuration the sample remains stationary while both
the X-ray source and the detector move by the same angle
Fig. 4.14 Typical spectra of diffracted X-rays
X-ray peaks at different angles of rotation in Fig. 4.14 originate from the X-rays
diffracted by different families of crystallographic planes in the crystal. The intensity
of the peak is proportional to the density of planes of a given family. As the crystal
is rotated, the angle of incidence θ changes for different families of crystal planes
and that family of planes for which Bragg’s law 2d sin θ = nλ gets satisfied diffracts
X-rays producing constructive interference and a maximum in intensity.
4.3 Dual Nature of Matter
The description of the motion of particles in classical physics got revolutionised

with the hypothesis of Louis Victor de Broglie (pronounced as de Broy) that every
radiation has a particle-like nature and vice versa. This was referred to as the dual
nature of matter. It meant that a moving particle of matter should exhibit wave-
like properties under suitable conditions. It was mentioned that since the Nature is
symmetrical the matter and energy should also be symmetrical. Thus, both radiation
as well as matter should not be different.
It is often questioned as to what led de Broglie to postulate that matter should
have wave-like properties and that waves must possess matter like aspects. Though no
clear answer may be given, but it may be said that by 1924 when de Broglie proposed
the dual nature of matter, Einstein’s mass–energy equivalence (1905), Bohr model
of hydrogen atom (1913) and the phenomena of Compton scattering of gamma rays
(1923) were already known. Einstein’s mass energy equivalence clearly indicated that
mass is a form of energy; similarly, it was postulated in Bohr’s model of hydrogen
atom that energy difference between two energy states of electron may be emitted
in the form of electromagnetic waves during the process of de-excitation of the
atom, showing that waves are also a form of energy. In other words, these two
postulates clearly established the equivalence of mass, energy and waves. The only
remaining question that wave may also behave like a particle was established by
the phenomena of Compton scattering, where a high-energy gamma ray behaving
like a particle strikes a stationary atomic electron, scatters the stationary electron
4.3 Dual Nature of Matter 211
in some direction, loses some energy and is itself scattered with reduced energy in
the complimentary direction. All these facts might have led de Broglie to postulate
matter waves.
According to de Broglie, the wavelength associated with a material particle is
given by, λ = h/ p, where h is Planck’s constant, while p = mv is the linear
momentum of the particle of mass m moving with velocity v. The validity of the
de Broglie relation can best be tested by the results of the experiment. It may be
remarked that the above equation is satisfied by a photon as well because the linear
momentum of a photon p = hν/c. Thus, hp = νc = λ.
De Broglie’s postulate essentially says that a particle of mass m, moving with
velocity v, has an associated wavelength λ = h/mv. In order to appreciate de Broglie
wavelength λ, it is required to understand what is meant by the term wavelength.
Wavelength is essentially the uncertainty in the position of a moving particle. The
location of a moving particle at any instant in its path of motion is uncertain by the
amount of its wavelength. It is obvious that if the physical dimension (the size) of
the moving object is smaller than the associated de Broglie wavelength, then only it
will be possible to carry out experiments to see the effect of wavelength. In case the
size of moving object is larger than the associated wavelength, it may not be possible
to experimentally detect the wavelength associated with the particle.
It can be seen from the expression for de Broglie wavelength that a particle of large
mass will exhibit a smaller wavelength and vice versa. This is the reason, why in our
daily life the wave-like character of even a cricket ball thrown with 80 miles/h is not
observed. In general de Broglie wavelength of only microscopic atomic and nuclear
particles moving with high velocities may be measured experimentally. According
to Einstein’s theory of relativity, the mass of a moving body depends on its speed of
motion. For particles moving with high velocities relativistic express for mass given
below must be used.
m0
m=/ (4.20)
v2
1− c2
Here, m0 is the rest mass of the body.

√ Since the momentum p and kinetic energy E are related by the expression, p =
2m E, as such the de Broglie wavelength may also be written as:
h
λ= √ . (4.21)
2m E
SAQ: What is the rest mass energy of an electron in MeV?

SAQ: What is the value of de Broglie wavelength for a stationary body and what
does it mean?
Solved Example SE4.1 Calculate the de Broglie wavelength of a ball of mass 100 g
moving with a speed of 30 m/s.
Solution: Given, m = 0.10 kg; v = 30 m/s
h h 6.64 × 10−34 J s
λ= = = = 4.21 × 10−34 m
p mv (0.10 kg) × 30 ms
As can be seen from the value of the wavelength obtained above that it is so
small that it is not even measurable with present day instruments. This shows the
reason why the wave nature of macroscopic objects is not observable in our daily
life. On the other hand, for microscopic particles the wave-like nature is significant
and observable.
As an example, let us consider that an electron of m and charge e is accelerated
across a potential of V volts, then the electron gains a kinetic energy K .E. = eV .
p2
√
Since K .E. = 2m , therefore, p = 2meV .
The de Broglie wavelength λ = √2meV h
.
If in the above expression we substitute the numerical values of Planck’s constant
h, mass of electron and charge of electron, then the above expression reduces to λ =
1.227
√
V
nm, where V is the accelerating potential in volts. For a value of V = 150 V,
the value of wavelength comes out to be λ = √ 1.227
150
nm = 14.247
1.227
nm = 0.100 nm. This
value is comparable to the order of spacing between the atomic planes in crystals. It
indicates that the particle nature of electrons could be verified by crystal diffraction
experiments similar to the X-ray diffraction. The experimental verification of the de
Broglie hypothesis has been described in detail in the next section. Louise Victor de
Broglie was awarded the 1929 Nobel Prize in Physics for his discovery of the wave
nature of electrons.
4.3.1 Davisson and Germer Experiment
Clinton Davisson and Lester Germer were involved in the study of the surface prop-
erties of Nickel samples using a beam of low-energy electrons at Bell laboratories,
USA, since 1923. However, the remark by Walter M. Elsasser (scientist at Gottingen,
Germany) that electron scattering by crystalline solids may be used to test the wave
nature of electrons, led Davisson and Germer in 1927 to repeat their experiment,
now with the view to look for the wave nature of electrons.
Experiment carried out by Davisson and Germer provides direct verification of
De Broglie hypothesis of the wave nature of moving bodies and demonstrated that
moving electron has an associated wave. The typical experimental setup used by
Davisson and Germer is given in Fig. 4.15. The thermionic emission of electrons
from a hot tungsten filament was used to provide a beam of electrons. These electrons
were accelerated by applying suitable potential difference with the help of a battery
as shown in Fig. 4.15.
The electron beam was collimated by allowing them to pass through a cylindrical
arrangement with a fine slit. The electron beam is incident on the Nickel crystal
having ordered arrangement of atoms, at normal incidence. Nickel atoms diffracted/
Fig. 4.15 Schematic diagram of Davisson–Germer experimental arrangement
scattered the incident electrons. In the experimental setup Nickel target crystal had an
arrangement of rotation at different angles in a plane. The intensity of the electrons
scattered by the crystal in a given direction was measured with the help of movable
detector. The whole experimental arrangement was placed in a highly evacuated
chamber. Several experiments were carried out and intensity of scattered beam at
different angles was recorded for different accelerating voltages. The observed curves
are plotted in polar coordinates in Fig. 4.16. Surprisingly, instead of a continuous
variation of scattered electron intensity with angle distinct maxima and minima were
observed whose position depended on the electron energy. It was observed that there
is a pronounced maximum that appear at ϑ = 50◦ angle of scattering (with respect to
the direction of the incident beam), when the accelerating potential is 54 V. Further,
increasing the accelerating potential indicated that the bump like maximum decreases
and becomes almost insignificant as the accelerating potential reaches to 68 V.
Figure 4.17(i) shows the variation of the intensity of 54 eV energy electrons with
angle of scattering/diffraction by the Ni-crystal. As may be observed in this figure
a prominent maximum in the observed intensity occurs at angle 50°. Figure 4.17(ii)
shows the direction of the incident beam of electrons which is normal to the top face
of the Ni-crystal. However, the incident beam makes an angle of 65° with family of
dominant crystal planes that are responsible for the diffraction of incident electrons.
Further, from X-ray diffraction experiments it was known that the separation between
successive crystal planes d is 0.91 Å.
The reason for the prominent bump like state at ϑ = 50◦ at accelerating potential
of 54 V may be understood as due to the diffraction of electron waves by the crystal
planes in the target Ni-sample. Figure 4.17(i) shows the electron diffraction pattern
for 54 eV electrons as a function of diffraction angle.
Fig. 4.16 Polar plots of the diffraction pattern of electrons by Nickel crystal for different
accelerating voltages
Fig. 4.17 (i) Diffraction pattern of 54 eV electrons by Ni-crystal. (ii) Crystal plane family and
diffraction geometry
The de Broglie wavelength associated with electrons accelerated through 54 V

may be calculated as,
1.227 1.227 1.227

λde Brg = √ nm = √ nm = = 0.168 nm = 1.68 Å.
V 54 7.348
One may also calculate the wavelength of the waves that have produced the
observed maximum using Bragg’s law. While applying Bragg’s law one must
remember that the angle of scattering that the incident electron beam makes with
the crystal plane is 65° and that the separation between planes is 0.91 Å,
λBragg = n2d sin θ

Substituting order of diffraction n = 1, θ = 65° and d = 0.91 Å in the above

equation, the wavelength obtained from Bragg’s law λBragg comes out to be
λBragg = 1 × 2 × 0.91 × sin 65 Å = 2 × 0.91 × 0.906 Å = 1.66 Å.
It may be observed that there is very good agreement between the de Broglie
wavelength λde Brg (= associated with 54 eV electrons = 1.68 Å) and the wavelength
λBragg (= 1.66 Å) obtained using Bragg’s law. It clearly proves that a beam of electrons
behaves both as a beam of material particles as well as a beam of waves of wavelength
given by de Broglie formula.
Davisson and Germer experiments do provide a direct verification of de Broglie
hypothesis of the wave nature of moving bodies. Soon after the publication of the
results from Davisson–Germer experiment many more and detailed experiments
were performed all of which confirmed the existence of matter waves. The electron
diffraction was studied by G. P. Thomson using X-ray powder diffraction method,
where an accelerated fine beam of electrons was made to hit normally on to the thin
metallic foil. The other side of the foil was exposed to a photographic plate. As the
electron beam passes through the foil, electrons of the incident beam get diffracted
by the grating like crystal structure in the foil forming a diffraction pattern on the
photographic plate. The diffraction pattern of bright co-centric circular rings around
a central spot got reviled when the photographic plated was developed. In order
to test that the diffraction pattern formed on the photographic plate is due to the
diffraction of electrons of the incident beam, a magnetic field was applied between
the source of electrons and the metallic foil. As expected, the diffraction pattern got
disturbed by the magnetic field confirming that the diffraction pattern is truly due
to the diffraction of de Broglie waves associated with accelerated electrons. Since
then many instruments particularly electron microscopes that use de Broglie waves
associated with accelerated electrons have been constructed and are in wide use.
Since the resolving power of a microscope depends on the wavelength of the waves
used in the instrument, de Broglie waves of very small wavelength may be produced
using high-energy electrons. De Broglie waves of very small wavelengths associated
with high-energy electrons are used in electron microscopes.
It might interest you to know that the postulate of matter waves proposed by de
Broglie was a part of his Ph.D. thesis.
Confirmation of waves associated with particles through several experiments
involving diffraction and interference of particle waves put de Broglie’s theory on
firm footing. However, two big questions regarding the matter waves also spring
up in the background of dual nature of matter. One big question was: what is the
velocity of the matter waves? And the second question was: what (material/field/
function) makes de Broglie waves? For example, in case of waves in a pond or sea,
it is water that moves up and down making the wave, in a wave on a string, different
segments of the string moves to generate wave and in electromagnetic waves (light)
it is the intensity of electric and magnetic fields that varies with time and generate the
electromagnetic wave; so the question is: time variation of which quantity constitutes
matter waves?
(a) Velocity of de Broglie wave
Let us address the first question: what is the speed of de Broglie (or matter) waves.
It may be recalled that in order to explain phenomenon like photoelectric effect and
Compton scattering it was necessary to assume that electromagnetic (EM) waves
(light and gamma rays) have an associated particle, called photon. The energy E phot
of photon, the momentum pphot of photon and the wavelength/frequency νphot of
electromagnetic wave are related according to the following relations:
h E phot hνphot c
E phot = hνphot ; λphot = ; pphot = = ; νphot = ,
pphot c c λphot (4.22)
and νphot λphot = c
Here, ‘c’ is the velocity of light in vacuum and ‘h’ is Planck’s constant.
De Broglie argued that if waves have a particle associated with them, then on the
basis of symmetry, moving material objects must also have associated waves, which
he called matter waves. De Broglie assumed that relations corresponding to the set
of relations given by Eq. (4.22) may also be written for matter waves.
Energy of the particle that carry matter waves E matt = hνmatt (4.23)
But from Einstein’s mass energy relation

m0
E matt = mc2 = / c2 (4.24)
v2
1− c2
Equating Eqs. (4.23) and (4.24) one gets,
mc2
νmatt = (4.25)
h
Also the wavelength λmatt of matter waves may be given as
h h
λmatt = =⎛ ⎞ (4.26)
mv
⎝ / 1
⎠m 0 v
2
1− vc2
In analogy to the expression for the velocity of EM wave c = λphot νphot , one may
write the velocity of matter waves as
h mc2 c2
Vmatt = λmatt · νmatt = · = (4.27)
mv h v
In Eq. (4.27), Vmatt represents the velocity of the matter wave, while ν is the
velocity of the particle of rest mass m0 . Now there is a contradiction. According to
the theory of relativity, no material particle can move with the velocity of light (c),
and therefore, the velocity v of the particle must be small than
the velocity of light,
c2
i.e. ν < c, and hence the velocity of the matter wave Vmatt = v may exceed the
velocity of light.
This puts a big question mark; what does this mean? The answer is that we have
to re-visit the dynamics of wave propagation.
In general, two velocities may be associated with a wave; they are: (i) the phase
velocity and (ii) the group velocity.
(i) Phase velocity
To understand the concept of phase velocity, let us consider a string or a wire held
fixed at two points in the x-direction. Suppose the wire is plugged in the y-direction
so that it vibrates in y-direction. The displacement ‘y’ of any point on wire at time
‘t’ may be given by,

x
y = A cos 2π ν t − (4.28)
Vphase
Here ν is the frequency and Vphase the wave speed, i.e. the speed with which the
wave travels down the wire. As a matter of fact Vphase is the speed with which the
displacement y travels along the wire in x-direction. Pick up any particular displace-
ment (in vertical direction) y1 at point X 1 of the wire at instant ‘t 1 ’. At the next instant
t 2 , the same displacement y1 will move to another location x 2 of the wire, and next at
instant t 3 , displacement y1 will reach at point x 3 on the wire. In this way the location
of displacement y1 is moving along the length of the wire in the x-direction. No part
of wire is moving in x-direction. Wire is moving (vibrating) in Y-direction. It may
thus be observed that no material particle moves along the x-direction, while the
location of a given displacement moves along the x-direction with speed Vphase , as
shown in Fig. 4.18. This speed, with which the phase of the wave travels, is called the
phase velocity of the wave. Since no material particle travels with the phase velocity,
the phase velocity may have a value equal or even greater than the velocity of light.
2
V matt = cv given by Eq. (4.27) represents the phase velocity of the matter
wave. Further, in matter waves no matter/particles are vibrating in y or in any
other direction.
(b) Group velocity
Equation (4.28) represents a progressive wave,i.e. a wave which is moving in x-
direction with velocity V phase . A negative sign in t − Vphase tells that the disturbance
x

is moving in +ve X-direction, while a positive sign in t + Vphase x
means that the
wave is moving in negative X-direction. It is often more useful to write Eq. (4.28) in
a slightly different form as given below,
Fig. 4.18 In case of a wire

stretched in x-direction and
vibrating in y-direction, the
location of a particular
vertical displacement moves
in x-direction with phase
velocity

x
y = A cos 2π ν t −
Vphase
Or

2π νx
y = A cos 2π νt − (4.29)
Vphase
But 2π ν = ω (the cyclic frequency of the wave) and Vphase = ν · λ.

Substituting these values in Eq. (4.29) gives

2π νx 2π x
y = A cos ωt − = A cos ωt − (4.30)
νλ λ
One defines the wave number (or wave vector) k = 2π λ

. With this substitution and
using the fact that cos(−θ ) = cos(θ ), Eq. (4.30) reduces to:
y = A cos[(kx − ωt)] (4.31)
where
2π mc2 2π m 0 c2
ω = 2π ν = = √ (4.32)
h h 1 − v2 /c2
And
2π 2π mν 2π m 0 v
k= = = √ (4.33)
λ h h 1 − v2 /c2
It follows from Eqs. (4.32) and (4.33) that:
ω 2π ν
= 2π = νλ = the velocity of the wave (4.34)
k λ
In three dimensions the wave equation becomes

y = A cos k→ · r→ − ωt , (4.35)
here k→ · r→ represents the dot product of vector r and vector k.

Equations (4.31) and (4.35) represent in one dimension and in three dimensions
travelling transverse waves which are polarised in y-direction.
The question now arises that if these ways are matter waves associated with
some material particle then where does the motion of the particle is hidden in these
equations? Clearly, these waves represent the motion of the phase of the wave which
travels with speed V phase . These waves do not represent the motion of any material
body.
Motion of a material body may be introduced in the wave equation only and only
by superimposing one or more than one waves. When there is more than one wave
with slightly different values of ω and k, travelling through the same region of space,
they undergo interference and generate a wave packet. It may be shown that the wave
packet produced by the superimposition of waves (or as a result of modulation of one
wave by other waves) may represent a material body and the speed with which this
wave packet travels is called the group velocity of the packet. The group velocity
Vgroup may be the velocity of the material particle associated with de Broglie waves.
Production of beats in acoustics, when two sound waves of slightly different
frequencies modulate each other, is well known. The pitch of the resultant sound
increases and decreases periodically. This phenomenon of the increase and decrease
of sound intensity is called the beat formation. When beads are formed by two sound
waves of slightly different frequencies, the number of beats or the repetition rate
of beads depends on the frequency difference of the two sounds. It may be noted
that though both sound waves individually travel with the same velocity (velocity of
sound in air), but the rate of beat formation depends on the frequency difference of
the two waves (Fig. 4.19).
To keep our discussion simple, let us consider two one-dimensional waves y1 and
y2 having same amplitudes A but slightly different cyclic frequencies ω and the wave
numbers k, moving in the same part of the space. The waves may be represented as:
Fig. 4.19 Superposition of

two waves of slightly
different frequencies
produces wave packets that
move with the group velocity
y1 = A cos(ωt − kx) and y2 = A cos[(ω + Δω)t − (k + Δk)x].
The two waves will interfere and the resultant wave may be represented as:
y1 + y2 = A cos(ωt − kx) + A cos[(ω + Δω)t − (k + Δk)x] (4.36)
Expression (4.36) may be simplified using the following relations:

α+β α−β
cos α + cos β = 2 cos cos and cos(−θ ) = cos θ
2 2

1 1
y1 + y2 = 2 A cos[(2ω + Δω)t − (2k + Δk)x)] cos (Δω)t − (Δk)x (4.37)
2 2
Since Δω and Δk are much small compared to 2ω and 2k, respectively, one may
take
2ω + Δω ≈ 2ω and 2k + Δk ≈ 2k.
With these approximations, Eq. (4.37) reduces to:

Δω Δk
y1 + y2 = 2 A cos(ωt − kx) cos t− x (4.38)
2 2
Equation (4.38) represents a wave of angular frequency ω and wave number k

that has been modulated by a wave of angular frequency Δω 2
and wave number Δk 2
.
The phase velocity of this modulated wave (see Eq. 4.34) is given by Vphase = ωk .

Term cos Δω 2
t − Δk
2
x in Eq. (4.38) represents a wave packet which travels with the
Δω
velocity Δki . The velocity with which the wave packet travels is called the group
Δω
velocity, and therefore, the group velocity Vgroup = Δki .
The magnitude of the group velocity may be obtained using Eqs. (4.32) and (4.33).
From Eq. (4.32)
2π m 0 c2
ω= √
h 1 − v2 /c2
Therefore,
dω 2π mv
= 3/2 (4.39)
dv 2
h 1 − vc2
Also from Eq. (4.33)
2π m 0 v
k= √
h 1 − v2 /c2
Differentiating k with respect to ν gives,
dk 2π m
= 3/2 (4.40)
dv 2
h 1 − vc2
It follows from Eqs. (4.39) and (4.40) that the group velocity
dω dω dk 2π mv 2π m
Vgroup = = / = 3/2
/ 3/2 = v (4.41)
dk dv dv 2 2
h 1 − vc2 h 1 − vc2
Thus, the group velocity (velocity with which the wave packet or the envelope
of the modulated wave travels) comes out to be equal to the velocity of the particle
associated with the matter waves, while the phase velocity of matter waves ωk = cv .
2
SAQ: Show that in a non-dispersive medium group velocity of an EM wave is equal

to the phase velocity.
(b) What makes matter waves
Waves are generated in space by the time variation/fluctuation of some material/

fields or some other quantity. For example, sound consists of pressure difference in
the medium, water waves are produced by the fluctuation in heights of water column,
and electromagnetic waves consists of fluctuating electric and magnetic fields. Matter
waves or de Broglie waves are produced by the time variation/fluctuation of the
wavefunction ψ of the associated material particle.
Schrodinger while developing the wave theory of particles (also called the
quantum mechanics) introduced the concept of wavefunction of a particle, denoted
by Greek letter ψ, which is a function that contains all properties of the particle. Once
the wavefunction for a system/particle is known, all properties of the system/particle

are known, and the system or the particle is completely defined. Wavefunction ψ is
generally complex and it cannot be directly measured. In general ψ is a function of
time and position (x, y, z).The probability of finding the object, for which ψ is the
wavefunction, at some point P (x, y, z) at time ‘t’ is proportional to the value of the
quantity ψ ∗ ψ at point P. When the value of ψ ∗ ψ at point P is 1, it means that the
object is there, and when the value is zero, it means that the object in not at point
P. It may be noted that the wavefunction of an object tells about the probability of
finding the object at some particular point in space and time, but it certainly does not
mean that the object has spread out in a wave. Further properties of wavefunctions
will be discussed while dealing with Schrodinger’s equation in quantum mechanics.
4.4 Some Examples of the Failures of Classical Approach

and Success of Quantum Approach
Classical physics is built on Newtonian mechanics, thermodynamics and Maxwell’s

theory of electromagnetism. Many theories of classical physics break down when
applied to microscopic systems or to objects moving with high velocities comparable
to the speed of light. An essential feature of classical approach is the assumption that
physical variables like energy, angular momentum, etc. vary continuously. As a matter
of fact the classical approach could not even explain the existence of an atom. Many
other processes observed experimentally could not be explained by the classical
theories, for example the photoelectric effect and Compton scattering. However, the
biggest problem was encountered in explaining the energy distribution of blackbody
radiations and the specific heat of solids. Failure of classical approach in explaining
the above mentioned processes will be discussed in the following.
4.4.1 Stability of the Atom and the Nature of Atomic Spectra
Rutherford in 1911 carried out some ingenuous and revealing experiments in which
thin gold foils were bombarded by energetic alpha particles. Scattering of incident
alpha particles by scattering angles as large as 180° established that there is a body
at the centre of each atom where total positive charge and more than 99% mass of
the atom are contained. The central body was named ‘Nucleus of the atom’, term
nucleus being borrowed from biology. Earlier, J. J. Thomson has already discovered
electron in 1897 and experiments with cathode-ray tube have proved that electron
is an essential constituent of all atoms. Soon after the discovery of atomic nucleus,
several theories for the structure of nuclear atom were proposed. The most convincing
model for atomic structure was the planetary model where it is assumed that electrons
in an atom revolve round the nucleus in circular orbits at different distances from the
4.4 Some Examples of the Failures of Classical Approach and Success … 223
nucleus, like planets revolve round the sun in solar system. The planetary model of the
atom was readily accepted because of its simplicity and compelling similarities with
the planetary system. For example, Coulomb force of attraction between the electron
and the nucleus may well be compared to the gravitational force of attraction between
the sun and the planet and that both these forces have almost similar dependence on
distance. Planetary model of the atom was also attractive as philosophically it was in
confirmation to the adage that big solar system is just an upscale of atomic system. At
the first sight planetary model of the atom appeared to be more stable than the solar
system, because in case of the solar system planets move in outer space where there is
matter of very low density, planets therefore experience a force of drag which reduces
the orbit of the planet and in long run every planet is expected to fall down into the
sun. No such drag was expected in case of the atomic electrons as the planetary
model assumes perfect vacuum around the nucleus where electrons revolve.
The stability of the planetary model of the atom which was based on classical
laws of physics (Newton’s laws of motion and Coulomb’s law) was questioned by
the classical theory of electromagnetism put forward by Maxwell (1862) in the form
of four equations. Maxwell’s theory says that a charge at rest has an electric field
around it which is strongly coupled to the charge; a charge moving with uniform
speed carry both the electric and the magnetic fields strongly attached to the charge
in uniform motion. However, if the charge is accelerated, a part of the electric and
magnetic fields which were strongly attached to the charge get detached and move
out in space with the velocity of light. Thus an accelerated electric charge according
to the classical theory of electromagnetism will radiate electromagnetic field and
loses energy. Although no force of drag is faced by the electrons in the atom, but
because electrons are assumed to have been moving in circular orbits, their direction
of motion changing at each instant, they are in accelerated motion and must radiate
energy. If so, the planetary atom which is based on classical physics is unstable from
the same classical approach. It is estimated that all electrons of an average atom will
spiral back into the nucleus within 10−8 s.
Further, the dying atom should emit electromagnetic (EM) waves of all frequencies
as electrons in different orbits will lose energy at different rates. In summery it may
be said that planetary atom which was the only possible model of the atom is (i)
unstable from the view point of classical physics and that (ii) while dyeing atom
must emit EM radiations of all frequencies.
Experimental facts, however, contradict both the above-mentioned predictions
of the classical approach. Atoms are in general stable, and when they emit EM
radiations, the radiations are not of all frequencies, and atoms on de-excitation emit
discrete EM radiations of fixed energies. As a matter of fact that the emission spectra
of atoms of each element of the periodic table is characteristic of the atom or the
element, it is like the signature or the thumb impression of the atom/element.
Several attempts, within the classical framework, have been made to explain
discrete atomic spectra by making different assumptions about the motion of electrons
in the atom, but none has been successful.
Figure 4.20 shows a representative line spectrum of sodium atom. As may be
observed in this figure, the spectrum has several lines in the ultraviolet region which
Fig. 4.20 Emission spectra

of sodium atom
lie in the invisible part of the EM spectrum; however, the two prominent lines called
D1 and D2 lines of yellow colour (wave lengths 589.0 and 589.6 nm) are characteristic
of sodium. It is the light from these lines which is used in sodium vapour lamps. There
are few other lines in the sodium spectrum that lie in the region of infrared and are
not shown in the figure.
Sodium lamp Sodium vapour lamps are frequently used for street lighting and can
be easily identified by their signature yellow light. Excited sodium atoms in visible
region emit yellow light of two wavelengths 589.0 and 589.6 nm. Sodium lamp has
the advantage that they are very efficient; almost 80% of the electrical energy given
to the lamp is converted into visible yellow light. Further, the lumen output of the
lamp does not drop with age and the light of yellow colour emitted by the lamp is
the colour to which human eye is most sensitive.
Problems associated with the classical planetary model of atom were satisfactorily
addressed by the quantum mechanical model of Schrodinger which is discussed in
Chap. 5 of the book. In quantum mechanical model it is shown that electrons in an
atom are placed in different energy states defined by principal quantum number n,
orbital quantum number l and magnetic quantum number m.
4.4.2 Photoelectric Effect
The discovery of photoelectric effect is a story in itself. It is difficult to give the credit
of discovering photoelectric effect to one person alone. First signatures of photoelec-
tric effect appeared in an experiment carried out by German scientist Heinrich Rudolf
Hertz in 1887 to identify electromagnetic waves. Electromagnetic waves, predicted
by Maxwell in 1865, were a hot topic at that time. Hertz in an experiment tried to
generate EM waves by producing a spark between two electrodes which were kept
at a small distance from each other and were maintained at high potential difference.
He observed that production of discharge becomes easier when the cathode was illu-
minated with ultraviolet light. He concluded that ultraviolet light falling on metallic
cathode emits some radiations from the cathode that ionises the gas between the elec-
trodes making it easier to conduct the spark. The next year, in 1888 another German
scientist Wilhelm Hallwachs repeated Hertz experiment with a simple geometry. He
took a clean circular plate of Zinc and mounted it on an insulating stand. The Zinc
plate was attached by a conducting wire to a gold leaf electroscope. The electro-
scope was then given negative charger. In normal conditions the electroscope lost
its negative charge very slowly. However, when the Zinc plate was illuminated with
ultraviolet light, the electroscope lost its charge very fast. On the other hand if the
electroscope was charged positively, there was no quick leakage of positive charge
even when the Zinc plate was illuminated with ultraviolet light.
The picture remained unclear till 1899 when Thomson conclusively proved
that ultraviolet light falling on metallic cathode/Zinc plate causes electrons to be
emitted from the target. The process of emission of electrons from metallic surfaces
when illuminated with light is called photoelectric effect and the electrons the
photoelectrons.
Philipp Eduard Anton Lenard, who earlier worked as assistant to Hertz, studied
in details the properties of electrons emitted from metallic bodies when illuminated
by ultraviolet and other lights.
A typical experimental arrangement to study the photoelectric effect is shown in
Fig. 4.21. Here, a photosensitive plate C is placed opposite to the metallic plate A
inside an evacuated glass tube. A potentiometer setup is used to apply desired value
of potential difference across plates C and A. The evacuated glass tube has a quartz
window through which EM radiations of desired frequency and intensity from source
S may be made to fall on the photosensitive plate C. Quartz crystal is transparent
to most parts of the EM spectrum, and therefore, quartz windows are used to allow
electromagnetic radiations of wide wavelength range to enter the tube without any
substantial absorption, and ordinary glass on the other hand absorbs most of the
ultraviolet part of spectrum. Often, plate C is also called the cathode and plate A as
anode as they are generally kept, respectively, at negative and positive potentials.
Fig. 4.21 Experimental setup to study the photoelectric effect

The working of the experimental setup may be understood as follows. When light
from source S, of a given frequency and intensity, is made to fall on the cathode plate
C, photoelectrons with some kinetic energy (or speed) are emitted from the plate.
Now if anode plate A is given a positive potential +V with respect to the cathode
C, emitted electrons are attracted by plate A and are collected on it. Photoelectrons
picked up by the anode plate A flow through the external circuit constituting the
photoelectric current. The current, which may be in the range of few milli- to few
microamperes depending on the frequency and the intensity of the incident EM
radiations, may be recorded by the ammeter in the circuit. On the other hand if the
polarity of the applied potential is reversed, i.e. plate C is given a positive potential
with respect to plate A, photoelectrons emitted from plate C will be repelled by
the potential at plate A, as a result the photoelectric current in the circuit will get
reduced. Magnitude of potential V and its polarity may be easily controlled by the
potentiometer and commutator combination shown in the figure. Under the reverse
voltage condition when plate C is at positive potential and plate A is at negative
potential, the magnitude of the current in external circuit will get reduced because of
the repulsion of photoelectrons. Experiments using electromagnetic (light) radiations
of different frequencies and intensities were carried out in which current through the
external circuit was recorded for different magnitudes and polarities of the potential
difference V between plates C and A. The main observations of these experiments
were that the magnitude of the measured photoelectric current depends on the (i)
material of surface emitting electrons, (ii) intensity of the radiations incident on plate
C, (iii) potential difference V between the plates, only when plate A is at a negative
potential with respect to C. These observations are point wise further elaborated in
the following.
(a) Dependence of photoelectric current on frequency of incident radiation
In experiments where monochromatic EM radiations of different intensities and
frequencies ν1 , ν2 , ν3 , . . ., etc. were made to illuminate the cathode C and photo-
electric currents in the external circuit were recorded, it was found that
(i) When plate A was at a positive potential +V with respect to plate C, the photo-
electric current changed only with the intensity of the incident radiations and
remained constant when EM radiations of different frequencies but of same
intensities were incident on plate C, for all values of positive potential +V.
Moreover, the photoelectric current was recorded immediately without any
time lag at the instant the radiations hit the cathode.
(ii) When EM radiations of same intensity but of different frequencies were incident
on the cathode, as mentioned earlier, it was found that photoelectric current
remained constant for all values of positive potential (+V ); however, the current
becomes zero when the frequency of the incident radiation was reduced to
some value ν 0 or below this value. The maximum frequency ν 0 at which no
photoelectric current passes through the circuit is called the threshold or cut-
off frequency. This indicates that no photoelectrons are emitted when EM
radiations of cut-off frequency ν 0 or of frequency lower than this are made to
shine the cathode, no matter what is the intensity or for how long the radiations
are made to hit the cathode C. The magnitude of the threshold frequency ν 0
has been found to be different for different metallic cathode surfaces.
(iii) In the case of reverse voltage setting, when plate A was kept at negative potential
(−V ) with respect to plate C, it was observed that the photoelectric current for
incident radiations of all intensities and frequencies decreased sharply with
the increase in the magnitude of the negative potential of plate A, becoming
zero for a the negative potential (−V s ). Negative potential (−V s ) where the
photoelectric current (for all intensities and frequencies of incident radiations)
becomes zero is called retarding potential or cut-off potential and has been
found to have different values for different metallic surfaces used as cathode
plate C. Further the cut-off potential does not depend on the intensity of the
incident radiations.
Let us now try to understand, within the framework of classical physics, the process
of photoelectric effect. Classically, EM radiations like all other waves carry energy
and are a mode of energy transfer. When EM radiations fall on a metallic plate, they
deposit energy in the plate at a certain rate, rate being proportional to the intensity
and the frequency of the radiations. Amount of energy deposited in the cathode plate
will be proportional to the time for which the plate is exposed to radiations and also
to the intensity and the frequency of the radiations. The cathode plate contains atoms
of the metal which have electrons that are bound to the bulk material of the plate
with some binding energy, say w. This w is called the work function of the metal
and is equal to the amount of energy required to take an electron out of the metal
surface. Obviously, the value of w depends on the metal used for cathode. According
to classical physics, emission of photoelectrons from the cathode plate will happen
only if the energy deposited by the incident radiations is at least equal or more than
the work function w of the material. If classical picture of photoelectric effect is
true, then there should be some time lag between the irradiation and recording of
photoelectric current, particularly for very low-intensity and low-frequency incident
EM radiations. Further, incident radiations of any frequency must be able to eject
photoelectrons, and low-frequency radiations, which deposit energy at a lower rate,
should be able to deposit the required energy w in a longer time of irradiation.
Therefore, radiations of all frequencies must be able to produce photoelectrons and
photoelectric current; the classical approach could not explain why radiations of
frequency less than the threshold frequency could not produce photoelectric current.
Photoelectric current is constituted by the number of photoelectrons collected
per unit time at anode A. On the other hand, the rate at which photoelectrons are
emitted from cathode plate will depend on the rate at which energy is deposited in
cathode plate by incident radiations. If the intensity of incident radiations is high,
more photoelectrons will be emitted per unit time, and hence there will be large
photoelectric current (as observed experimentally). However, according to classical
picture, the amount of energy deposited in cathode will increase with time, and
therefore, the photoelectric current should not remain constant but should increase
with time. This contradicts the experimental observations.
Summing up it may be said that instantaneous emission of photoelectrons, occur-

rence of threshold frequency, dependence of photoelectric current only on the inten-
sity of incident radiations are some major issues that could not be addressed by the
classical approach.
(b) Dependence of photoelectric current on the intensity of incident radiation
In experiments where the intensity (I) of incident radiations of fixed frequency was
made to hit cathode C and the voltage between plates C and A was kept constant,
it was found that the current in the external circuit was directly proportional to the
intensity of the incident radiations as shown by graph of Fig. 4.22.
Let us analyse if linear relationship between the intensity of incident EM radiations
and the photoelectric current may be explained by the classical approach. Classically,
the only requirement for the emission of a photoelectron is that an amount of energy
w (or larger than that) is deposited on the electron so that it may overcome its
binding with the cathode. Since EM waves of higher frequencies deposit energy
in the material of plate C at a faster rate, more photoelectrons should have been
emitted by incident EM waves of higher frequencies. Since photoelectric current is
proportional to the number of photoelectrons emitted per unit time from the plate; the
photoelectric current should have been proportional to the frequency of the incident
radiation as well as the intensity of the incident wave. In actual experiment no linear
relationship between the frequency of the incident radiations and the magnitude of
the photoelectric current has been observed. In the light of the above discussion,
results of these experiments could not be explained on the basis of classical physics.
(c) Dependence of photoelectric current on the potential difference across the
two plates
In some experiments, cathode plate C was illuminated with monoenergetic EM radi-
ations of constant intensity and frequency ν (> ν0 ). The photoelectric current (I) in
the external circuit was recorded for different settings of voltages between plates C
and A. Graphs for two different values of intensities I 1 and I 2 (I 1 > I 2 ) of incident
radiations of a given frequency are shown in Fig. 4.23. The potential difference is
taken positive when plate A is at a higher potential than plate C. It may be observed
in the figure that when potential difference is positive, the magnitude of current
remains constant. In the case of reverse potential when A was negative with respect
Fig. 4.22 Variation of

photoelectric current with
the intensity of the
monoenergetic (fixed
frequency) incident
radiations
to C, current decreases (for both intensities) and ultimately becomes zero at the same
value of reverse (or retarding) potential −V s for all values of intensities of incident
radiations.
The above experimental observations may be explained by assuming that the
incident EM radiations of frequency ν (> ν0 ) falling on the surface of plate C emit
photoelectrons of all kinetic energies from zero up to a maximum value E max , where
E max depends on the frequency of the incident radiations. Further, the number of
photoelectrons emitted per unit time is proportional to the intensity of the incident
radiations. In case when plate A is at a positive potential +V with respect to C, all
photoelectrons (of all energies from 0 to E max ) emitted per unit time are collected by
plate A, irrespective of the magnitude of potential V, and the photocurrent remains
constant for all positive values of V. However in case of the reverse potential when
plate A is at a negative potential −V n with respect to plate C, photoelectrons emitted
from plate C get repelled by plate A; at small negative values of V n only some low-
energy photoelectrons are not able to reach A reducing the magnitude of photocurrent,
but when negative potential increases, more energetic photoelectrons are also not able
to reach plate A. Ultimately when V n attains the value −V s , even the most energetic
electrons of energy E max are also cut off and photocurrent becomes zero. For V n =
−V s , one may write:
1
E max = 2
m e vmax = eVs . (4.42)
2
Here m e , e and vmax are respectively the mass, charge and maximum velocity of
emitted photoelectrons.
The above experimental observation that photoelectrons emitted by a certain EM
radiation of frequency ν will have a maximum kinetic energy is not supported by clas-
sical approach. According to classical approach, EM radiations will keep depositing
energy in the target material, and therefore, energy available to photoelectrons in
Fig. 4.23 Dependence of

photoelectric current on the
potential difference between
plates C and A
excess to their work function should also go on increasing with the length of time
the EM wave is kept shining on the plate C.
(d) Dependence of photoelectric current on the frequency of incident light and
on the stopping potential
Figure 4.24 shows the results of an experiment in which EM waves of three different
frequencies ν1 , ν2 and ν3 with (v3 > ν2 > ν1 > v0 the threshold frequency) having
same intensities were incident on plate C, one at a time, and the photoelectric current
for both positive and negative voltage settings between plates C and A was recorded.
Since the magnitude of the photocurrent for any positive voltage on plate A
depends only on the intensity of the incident EM radiations, the straight-line curves
for the three frequencies overlap on each other and are shown by a single straight
line marked as saturation current. However, in case of retarding potential when
plate A is at negative potential with respect to plate C, curves for EM radiations of
three different frequencies cut the X-axis at different cut-off (or retarding) potentials
−V s1 , −V s2 and −V s3 . Since |V s3 | > |V s2 | > |V s1 |, the maximum kinetic energy E max
v3
of photoelectrons emitted by the EM wave of frequency ν 3 is highest. This shows

that the maximum kinetic energy of photoelectrons depends on the frequency of the
incident EM wave and increases with its frequency.
It may be noted that retarding or cut-off potential, (−V s3 ), for example, is a
v3
measure of the maximum kinetic energy E max of photoelectrons emitted by EM
radiations of a given frequency ν 3 . If one draws a graph between the maximum
kinetic energy of photoelectrons and the frequency of the incident EM waves for
a given material (metal) on plate C, a straight-line graph marked M 1 in Fig. 4.25
is obtained. Similar graphs for other metals M 2 and M 3 are also parallel straight
lines, but the straight lines for different target metals cut the x-axis at different points
marked as ν0M1 , ν0M2 , ν0M3 , etc. It means that the threshold frequencies, the minimum
Fig. 4.24 Effect of variation of frequency of incident light on stopping potential

Fig. 4.25 Graph showing

the dependence of maximum
kinetic energy of
photoelectrons on the
frequency of the incident EM
wave for different target
metals
frequency of the incident radiations below which no emission of photoelectrons takes

place, have different values ν0M1 , ν0M2 , ν0M3 , etc. for different target metals.
As already mentioned, classical theory of EM waves could not explain why waves
with frequencies below the observed threshold frequency (ν 0 ) could not emit photo-
electrons from the target metal. Classical approach also fails to explain why different
metallic surfaces have different values of threshold frequencies.
In the conclusion it may be said that classical electromagnetic theory put forward
by Maxwell could not explain the following experimental observations on photoelec-
tric effect: (a) no time lag between the emission of photoelectron and the striking of
EM radiations on target metallic plate. (b) Independence of the maximum kinetic
energy of photoelectrons from the intensity of impinging EM radiations. (c) The
presence of a cut-off or threshold frequency, EM radiations of frequencies below
the threshold frequency cannot emit photoelectrons no matter for how long the
target material is irradiated by EM radiations.
(e) Dependence of cut-off (threshold) frequency on the type of cathode surface
Einstein’s quantum mechanical approach that explained and removed all the above-
mentioned drawbacks of classical theory will be discussed in the following.
4.4.3 Quantum Theory of Photoelectric Effect
According to the classical picture of EM radiations put forward by Maxwell, radi-

ations are electromagnetic waves, the energy contents of which are proportional to
the square of the amplitude of the wave. However, this description of EM waves
could not explain various experimental observations on the process of photoelectric
effect as mentioned in the last section. Similar discrepancies were also encountered
by Planck and others in explaining the energy distribution of EM radiations emitted
by a blackbody. It was Planck who in 1900 made the bold proposition that energy is
absorbed and emitted not continuously but in small energy packets which he called
energy quanta. Einstein in 1905 proposed the quantum theory for photoelectric effect,
borrowing the idea of Planck that energy is emitted or absorbed in energy packets
or energy quanta. Einstein assumed that EM radiations are made up of tiny energy
packets, called photons which move with the velocity of light in vacuum. The energy
ν
E pho of a photon of light of frequency ν is given by,
ν
E pho = hν (4.43)
where h is Planck’s constant.

Therefore, a beam of monochromatic light (EM radiations) of frequency ν may be
considered to be a bundle of photons, each of energy hν, moving with the velocity c,
the velocity of light in the medium. The intensity of the EM wave is proportional to
the number of photons in the bundle. It may be mentioned that photon is not a material
body or particle, in other words the rest mass of photon is zero, and it existsEonly
ν
ν ν
in motion and possesses energy (E pho = hν) and momentum ppho = hνc
= pho
c
.
A photon when interacts with a material particle, like an electron, obeys the laws of
conservation of energy and momentum. Photon aspect of EM wave is the counterpart
of its wave aspect and is in confirmation with principle of duality. It may also be
mentioned that the wave and the particle aspects of EM radiations do not show up
simultaneously; when radiations show wave aspect (in interference and diffraction),
it does not show photon or particle aspect, and similarly when radiations exhibit
particle aspect (photoelectric effect, Compton scattering, etc.), they do not show
wave aspect.
Einstein explained the photoelectric effect assuming that when a metallic surface
is illuminated with EM radiations of frequency ν, an incoming photon of EM radiation
of energy hν hits an electron in the atom of the target metal. If the energy (hν) of the
photon is more than the binding energy (w) of the electron with the bulk metal, the
incident photon may be absorbed and the struck electron gets ejected from the atom
as a photoelectron. The maximum kinetic energy E max of the emitted photoelectron
will be very nearly equal to the difference (hν − w) of the photon energy and the
electron binding energy in bulk metal. The electron binding energy with the bulk
metal w is called the work function of the metal, which has different values for
different metal surfaces.
Figure 4.26 shows that a photon of an incident EM radiations of frequency ν,
energy E pho (= hν) and linear momentum ppho = hν c
interacts with an electron of
the K-shell of the target atom. The photon imparts its total energy and the momentum
to the target atom and vanishes. In case the energy deposited by the photon is more
than the work function w of the electron, the electron may be ejected from the target
metallic surface with some kinetic energy denoted by E elec and linear momentum
pelec . To apply the laws of conservation of energy and momentum to the interaction,
one may look to the interacting system before the interaction and after the interaction.
There were two entities, the photon and the target atom, before interaction, and there
are two entities, the photoelectron and the residual atom, after interaction. Energy
and momentum need to be conserved between these entities. If it is assumed that the
target atom is at rest, the linear momentum pumped by the photon must be shared
between the photoelectron and the residual atom. As such, the residual atom must
recoil in a particular direction to conserve the input linear momentum. Some energy,
say E reco , is consumed in this recoil. It follows from the conservation of energy that:

E elec = E pho − w − E reco (4.44)
Since E reco may have different values for different photoelectrons, the energy
of emitted photoelectrons may differ from each other by the amount oi the recoil
energy, which is very small. Another reason for the difference in the kinetic energies
of photoelectrons is the depth of the atom (which has lost the photoelectron) from the
front surface. When photoelectron is generated deep inside the metallic target, it may
lose some of its kinetic energy in reaching the surface. Thus, differences in the value of
recoil energies and in the energy loss while coming out of the metallic surface produce
distribution in the kinetic energy of emitted photoelectrons. Maximum kinetic energy
is possessed by the photoelectron which is produced just at the front surface of the
target and for which recoil energy is a minimum.
In Fig. 4.26 it is shown that the photoelectric effect is taking place with an electron
of the K-shell. It is because the probability of photoelectric effect with electrons of
inner shells, like the K- or the L-shells of the atom, is a maximum. The reason for
this is the fact that in photoelectric effect conservation of linear momentum demands
that the residual atom (left after the emission of photoelectron) must recoil. This may
happen only when the emitted photoelectron was tightly bound with the atom and
may easily transfer the excess linear momentum to the residual atom. Since K- and
L-shell electrons are most tightly bound with the atom, photoelectric effect is more
likely to take place with these electrons.
Einstein’s quantum mechanical model of photoelectric effect may explain simul-
taneous emission of photoelectrons, without any time lag, with the irradiation of
the target metal surface by EM radiations, if the energy of the incident photon is
more than the work function w of the target metal. A rough estimate of the threshold
frequency ν 0 may be made from the work function w of the target metal as:
Fig. 4.26 Emission of photoelectrons by the absorption of incident photon

w
v0 ≈ (4.45a)
h
Since different metallic surfaces have different values of the work function, the
cut-off or threshold frequency has different value for different materials.
The intensity of incident EM radiations, according to the quantum approach, is
proportional to the number density (number per unit volume) of photons in the beam.
The number of photoelectrons emitted per unit time is also proportional to the number
density of photons in the incident bam, which in turn is proportional to the intensity
of the beam. Therefore, the photoelectric current which is constituted by the emitted
photoelectrons is also proportional to the incident beam intensity as shown in graph
of Fig. 4.23.
It may be remarked that the quantum approach to photoelectric effect given
by Einstein has been able to explain all experimentally observed facts about
photoelectric effect and could address the anomalies posed by the classical approach.
4.4.4 Work Function
In general one talks about the ionisation energy of an atom, which is the energy
required to take out an electron from the atom when it is in gaseous state. Situation
changes when one considers atoms in a solid, particularly in case of metallic solids.
The structure of a metallic solid may be described in terms of a positive ion lattice
surrounded by a cloud of de-localised electrons. Since in metals, an electron of the
cloud is not bound with an individual atom, the concept of ionisation energy is not
applicable. Instead the concept of work function is used; work function (w) may be
defined as the energy required in taking out one electron (of the electron cloud) to the
surface of the bulk material. In general for metals the work function (w) is smaller
than the ionisation energy (E ioniz ) of the corresponding atom. For example, in case
of copper w = 3.76 eV and E ionz = 7.52 eV and for Silver w = 4.34 eV and E ionz =
8.68 eV.
4.4.5 Residual Atom after the Emission of Photoelectron
The residual atom left after the emission of photoelectron is still excited and has an
electron vacancy in one of the inner most shells, like in K- or L-shells. Electrons
from higher shell, like M, N, …, etc., may fill the vacancy of the inner shell. This
transfer of electron from the higher shell to the lower shell is accompanied with the
emission of characteristic lines of the emission spectrum of the atom. For example,
if photoelectron is ejected from K-shell, then an electron say from the M-shell may
move to K-shell and quench the vacancy there, emitting K β X-rays of the atom.
Sometimes it may happen that electron from the M-shell goes to K-shell, but no
4.5 Blackbody Radiations and Their Energy Distribution 235
X-ray is emitted, instate excess energy (that might have gone out in the form of K β
X-ray) is given directly to some outer shell electron (which is loosely bound) and
that electron goes out of the atom. Electrons emitted in this way by the direct transfer
of excess energy are called Auger electrons and the process Auger effect.
Dependence of photoelectric effect on atomic number and energy of photon
The probability of photoelectric effect depends on the energy of incident radiation
and also on the atomic number Z of the target atom and may be represented by the
empirical relation;
Z 4.5
pPhotoelectric ∝ (4.45b)
hν 7/2
It follows from the above expression that the probability of photoelectric effect
is more for atoms of higher atomic number Z (heavier materials) and decreases with
the increasing energy of the incident photon.
SAQ: Why the photoelectric effect is said to be a bound state phenomenon?
4.5 Blackbody Radiations and Their Energy Distribution
Since in this section we will be dealing with thermal radiations, it will be appropriate
to define thermal radiations. It is a common observation that metallic objects when
heated emit electromagnetic radiations in the visible region and that the colour of
the emitted radiations changes with the temperature of the body. In general, not only
metals but all bodies emit EM radiations when they are at a temperature above the
absolute zero. The emitted EM radiations contain waves of many frequencies (or wave
lengths) the distribution of which depends both on the temperature and the material
of the body. The electromagnetic radiations emitted on account of the temperature
of any object are called thermal radiations. Thermal radiations emitted by a perfect
blackbody are termed as blackbody radiations. Thermal radiations emitted by an
object because of its temperature are quite different from the emission line spectra
of atoms or band spectra emitted by excited molecules.
The concept of blackbody and blackbody radiations in thermodynamics has orig-
inated from Kirchhoff’s law of thermal emission, given by German scientist Gustav
Robert Kirchhoff in 1862. There are several ways to express Kirchhoff’s law of
thermal emission. The original law which was in German may be translated in simple
English as: ‘if there is a region of space surrounded on all sides by perfectly insu-
lating boundaries so that no part of thermal radiations may leak through them and
if each part of the boundary is at the same constant temperature T, then the space
surrounded by the boundary is filled with thermal radiations which are characteristic
of temperature T alone’. He named these radiations as blackbody radiations at
temperature T and called the space surrounded by the boundary as the blackbody.
The characteristics of blackbody radiations at temperature T will be same as that of
thermal radiations emitted by a lump of lamp black held at constant temperature T,

hence the name blackbody radiations. Needless to say that blackbody radiations are
electromagnetic waves having a spectrum of frequencies or wavelengths distributed
in definite proportions. The condition that each part of the boundary should be at the
same constant temperature T ensures thermal equilibrium in the space bounded by
the boundary or in the blackbody.
Is it possible to make a real blackbody and to take out blackbody radiations
from it for studying their properties? The answer is big NO, because once a part
of the blackbody radiations are taken out of the cavity, thermal equilibrium will
get disturbed. Wilhelm Wien (German scientist), however, suggested that thermal
radiations very close (in character) to blackbody radiations may be obtained from
a very small hole made in the wall of a large container of any material and of any
shape, provided that the interior of the container (called thermal cavity) is kept at a
constant temperature T (i.e. in thermal equilibrium).
Since it is easy to make a source which may deliver thermal radiations very close
in character to blackbody radiations, many and very precise studies of blackbody
radiations at different temperatures have been made. A typical spectral distribution
of energy density of blackbody radiations for three different absolute temperatures
(4000, 5770 and 7000 K) are shown in Fig. 4.27.
A close look of the observed spectral energy density distribution reveals that
(i) at a given temperature, the energy is not distributed uniformly in the blackbody
spectrum. (ii) As the temperature of the blackbody increases, the energy density of all
wavelengths in the spectrum increases. (iii) The high frequency or short wavelength
emission cutoff shifts towards left side, towards shorter wavelength. (iv) Also, the
peak of the distribution plot, corresponding to the wavelength (or frequency) of
maximum emission (λmax ), shifts towards the short wavelength side and gets a bit
narrower. (v) The intensity of the radiation is found to change with the wavelength,
Fig. 4.27 Energy distribution of blackbody radiations at three different temperatures

and it increases exponentially at a faster rate and decreases exponentially with a

slower rate with the increase in wavelength.
4.5.1 Wien’s Displacement Law
After careful analysis of blackbody thermal radiation spectra at many temperatures,

Wien (1894–1896) put forward a relation between the wavelengths λmax and the
absolute temperature T. Since the law was derived purely on the basis of experi-
mental data of blackbody spectra without any theoretical backing, the law is empir-
ical and is called Wien’s displacement law. Mathematically, the empirical Wien’s
displacement law may be written in the following two forms:
λmax · T = Constant = 2.899 × 10−3 m K (4.46a)
Or
νmax
= Constant = 5.879 × 1010 Hz/K (4.46b)
T
It may appear surprising to note that the second form of the law given by
Eq. (4.46b) cannot be obtained by substituting νmax = λmax c
in Eq. (4.46a). The
reason is that the wavelength of maximum emission λmax is not a single wavelength,
but wavelengths lying in the range λmax and (λmax + dλmax ) are all wavelengths of
maximum emission. Similarly, frequency of maximum emission νmax is also not a
single frequency, but all frequencies in the range νmax and (νmax + dνmax ) are all
frequencies of maximum emission. Now, dλmax /= dνmax that means that the wave-
length and the frequency do not change at the same rate; hence, the two constants
are different.
Wien also proposed a law that may give the energy distribution in blackbody
spectrum. Unlike the displacement law, Wien derived his distribution law using the
laws of thermodynamics and Maxwell Boltzmann distribution law for the speed
of gas molecules. Essentially Wien used the concept of adiabatic compression of
blackbody radiations contained in an enclosure to reach a higher temperature. In
his derivation Wien made many approximations and his original derivation is quite
involved and lengthy. Also, in the light of the quantum theory of radiations it is not
of much relevance now. Skipping the derivation, Wien’s distribution law may be
given as,
C −D/λT
E λ dλ = e dλ. (4.47)
λ5
Here, E λ denotes the energy density contained in spectral range λ and (λ + dλ) of
blackbody spectra. C and D are two constants their values for a given temperature T
may be obtained by fitting the experimental spectrum at temperature T as such these
constants are temperature dependent. Though Expression (4.47) was derived using
the laws of classical thermodynamics and Maxwell Boltzmann distribution, but the
values of constants C and D need to be determined from experimental data, and the
expression is, therefore, semi-empirical.
The total emissive power E of the blackbody at temperature T, which may be
defined as the total energy radiated per unit time, may be obtained by integrating
Eq. (4.47) between the limits λ = 0 to λ = ∞, i.e.
∫∞ ∫∞
C −D/λT
E= E λ dλ = e dλ = σ λ4 (4.48)
λ5
0 0
where σ is a constant at a given temperature T and depends on the values of constants

C and D.
Equation (4.48) is nothing but mathematical representation of well-known
Stephan–Boltzmann law of thermodynamics. It may be observed that Stephan’s law
that was quite well-established law could be derived from Wien’s distribution law.
Similarly, it is possible to derive displacement law of Wien from his distribution law.
Successful derivation of these two laws gave good support to Wien’s distribution
law.
4.5.2 Failure of Wien’s Distribution Law
Wien’s distribution law may explain the shape of observed blackbody spectrum
at a given temperature, only qualitatively. The term λC5 e−D/λT of the distribution
formula may be considered to have two parts (a) λC5 and (b) e−D/λT . For small values
of λ exponential part (b) becomes large and over rides the effect of part (a); as
a result for short wavelengths, the energy density rises almost exponentially. On
longer wavelength side the exponential part (b) becomes very small and the fall in
the energy goes almost as λ−5 .
However several attempts to reproduce quantitatively the experimental energy
distribution curve of blackbody radiations at a given temperature T for the whole
range of wavelengths using Wien’s distribution formula failed; Wien’s distribution
law reproduced the lower wavelength part of the experimental energy distribution
but failed to reproduce the longer wavelength part of experimental distribution curve.
Moreover, assuming a nonzero value for the wavelength λ, if temperature T is set to
the value of infinity (∞) in Wien’s distribution
∫∞ formula, it is observed that the total
energy emitted by the blackbody E = 0 E λ dλ at an infinite temperature remains
finite. This is physically unjustified.
In short, Wien’s distribution formula failed as it has two drawbacks: (i) could
not reproduce the longer wavelength part of the experimental blackbody radiation
distribution and (ii) it predicts a finite value of energy being radiated by a blackbody
Fig. 4.28 Comparison of

experimental blackbody
spectra with predictions of
Wien and Rayleigh–Jeans
distributions
even at infinite temperature. Comparison of blackbody energy distributions obtained

experimentally and predicted by Wien’s formula is shown in Fig. 4.28.
4.5.3 Rayleigh and Jean’s Distribution Law
Strutt John William, Third Baron Rayleigh, a British mathematician better known
as Lord Rayleigh and Sir James Jeans in 1905 proposed another energy distribution
law known as Rayleigh–Jeans law to describe the energy distribution of blackbody
radiations. They argued that a blackbody cavity in thermal equilibrium at temper-
ature T (K) may be considered as if it is a cubical enclosure filled with standing
EM waves of different frequencies. They assumed that the walls of the blackbody
cavity contain some hypothetical oscillators that emit and absorb EM radiations
of different frequencies; the waves emitted by a particular oscillator interfere with
waves impinging on the oscillator producing standing waves. Putting the condition
that standing waves must have nodes at container walls, calculated the number of
modes of vibrations per unit volume dN in frequency range ν to (ν + dν) as:
8π ν 2
dN = dν. (4.49)
c3
In their derivation Rayleigh and Jeans assumed that the frequency ν of oscillators
varies continuously. Next they calculated the average energy per mode of vibration
E ave using the law of equipartition of energy of thermodynamics. The law says that
each degree of freedom has energy 21 kB T (kB being Boltzmann constant), and since
there could be two degrees of freedom: one corresponding to the kinetic energy and
the other of potential energy, the average energy for each vibrating oscillator becomes
E ave = kB T (4.50)
As such the average energy density of each oscillator of frequency in the range ν
to (ν + dν) may be written as:
8π ν 2 8π ν 2 kB T
E(ν) = E ave = (4.51)
c3 c3
Writing expression (4.51) in terms of wavelength, Rayleigh and Jeans obtained
the following formula for energy density of blackbody radiations as:
8π kB T
E λ dλ = dλ (4.52)
λ4
In Eq. (4.49) kB is Boltzmann constant and T the absolute temperature. It is worth
noting that Rayleigh–Jeans law in comparison to Wien’s distribution law does not
involve any new/unknown constants.
4.5.4 Failure of Rayleigh–Jeans Distribution
According to Rayleigh–Jean’s formula, the energy density in blackbody radiations

should always increase with the decrease in the value of the wavelength, which is
contrary to the observed experimental spectra where the energy density decreases
both for very short and for very large wavelengths.
The big flaw of Rayleigh–Jean distribution formula is that for shorter wavelengths
(high frequencies) the energy tends to become infinite. This is called ultraviolet
catastrophe.
It has been observed that Rayleigh–Jean distribution function can reproduce the
long wavelength component (after the peak) of the experimental blackbody spec-
trum, while the Wien’s distribution formula may correctly reproduce the shorter
wavelength component, below the peak, of the experimental spectra. Both of the
distribution formulae derived on the basis of classical physics (thermodynamics and
Maxwell Boltzmann distribution) fail to explain the blackbody spectrum in full.
Figure 4.28 shows the experimental energy distribution of blackbody radiations along
with predictions of Wien and Rayleigh–Jean.
4.6 Quantum Theory of Blackbody Radiations
Karl Ernest Ludwig Marx Planck better known as Marx Planck in 1900 put forward
the quantum theory for the energy density distribution of blackbody radiations. Like
Rayleigh–Jean, he also assumed that the blackbody cavity is filled with electromag-
netic waves which are continuously emitted and absorbed by some sort of oscillators
giving rise to the formation of standing waves. Using the condition that these standing
4.6 Quantum Theory of Blackbody Radiations 241
electromagnetic waves must have nodes at the boundaries of the cubical enclosure,
like Rayleigh and Jeans, Planck also obtained the same expression for the average
energy density of each oscillator of frequency in the range ν to (ν + dν) which may
be written as:
8π ν 2
E(ν) = E ave (4.53)
c3
At this stage Planck made a drastic assumption that oscillator cannot have contin-
uously variable energies; he said that oscillators may have only energies in integer
multiples of the quantity hν, where h is Planck’s constant. This assumption means
that there may be oscillators of energies, hν, 2hν, 3hν . . . nhν, where n is a positive
integer. Oscillators with energies 21 hν or 34 hν etc. were not possible. Next he calcu-
lated the probability p(n) of the mode with energy E n = nhν in thermal equilibrium,
using Boltzmann distribution law,

− kEnT
e B
p(n) = (4.54)
∑∞ − kEnT
n=0 e
B
The average energy density for mode of frequency ν is therefore,

∞ ∑∞ − kEnT
nhνe B
E ave = E n p(n) = n=0
(4.55)
∑∞ − kEnT
n=0
n=0 e B
In order to solve the above expression, we may substitute x = exp(−E/kT ).

Thus, the denominator of the above expression becomes,

n=∞
exp(−E n )/kT = 1 + x + x 2 + x 3 + · · ·
n=0
This is geometrical progression, and its addition will be given by:
1
(1/(1 − x)) = (4.56)
(1 − exp(−E/kT ))
Also,

n=∞

n E exp(−n E/kT ) = E x + 2x 2 + 3x 3 + 4x 4 + · · · (4.57)
n=0
d
xE 1 + x + x2 + x3 + · · ·
dx
E
d 1 Ex x exp kT
xE = = E 2 (4.58)
dx 1−x (1 − x)2 1 − exp kT
So the average energy of oscillators:
E exp(−E/kT ) E
E ave = = (4.59)
{1 − exp(−E/kT )} {exp(E/kT ) − 1}
This may be simplified into,
hν
E ave = hν . (4.60)
e kB T − 1
Substituting this value of average energy in Eq. (4.53), one gets,
8π ν 2 8π hv3 1
E(ν) = E ave = hν (4.61)
c3 c3 e kB T − 1
This is the Planck distribution function which reproduces the energy density
distribution of blackbody spectrum for all frequency or wavelength regions.
At the time Planck proposed his radical hypothesis, many scientists could not
believe mainly because Planck could not explain why the energies should be quan-
tized. Initially, his hypothesis explained only the experimental data on blackbody
radiation. It was mentioned that if quantization was observed for a large number of
different phenomena, then quantization would become a law. It was also remarked
that one needs to develop a theory that might explain that law. As things worked out,
Planck’s hypothesis was the starting point from which the modern physics grew and
developed.
Planck’s theory of blackbody radiations assumes that electromagnetic radiations
in the blackbody enclosure may have only discrete energies and the oscillators could
only lose or gain energy in the form of packets, referred to as quanta, of size hν, for
a given oscillator of frequency ν. The energy quanta of electromagnetic radiations
are called photons.
4.7 Compton Scattering of Gamma Rays
Electromagnetic radiations are produced when charged bodies are either accelerated
or decelerated (for example, the emission of continuous X-rays) and also when
electrons shift from one shell of the atom to another shell (example, characteristic
line spectra). In the case of the atomic line spectra, the energy of emitted photons
4.7 Compton Scattering of Gamma Rays 243
depends on the difference of energy of atomic levels or shells. Photons emitted in

atomic transitions may have energies of few ten of keV only.
Atomic nucleus contains neutrons and protons which are also distributed in energy
levels, as electrons are in an atom. Transitions of neutrons or of protons from a level
of higher energy to a level of lower energy in the nucleus of an atom give rise to
electromagnetic radiations of high energies, in the range of MeV. The electromagnetic
radiations that have their origin in the nucleus are called gamma rays and are often
denoted/represented by symbol γ . Gamma ray energies may lie in the range of few
keV to few MeV.
Gamma rays when incident on an atom are very unlikely to produce photoelectric
effect, as gamma ray photons have considerable energy and probability of photo-
electric effect decreases sharply with the energy of incident photon. Moreover, the
photoelectric effect could be explained only assuming a particle nature (photon) of
light.
American physicist Arthur Holly Compton at Washington University in 1923
discovered a phenomena in which high-energy EM radiations when incident on
an atom ejected an electron from the atom (like photoelectric effect) and a lower
energy EM radiation (unlike photoelectric effect) also appeared in the output. This
process was different from photoelectric effect as in photoelectric effect the incident
photon is lost and no photon appears in the output; there are photoelectron and
recoiling residual atom. Further, EM radiations’ characteristic of the incident atom
is also emitted following photoelectric effect. The process, in which high-energy
incident photon ejects an electron from the target atom and itself is scattered with
reduced energy, is called Compton scattering. Compton scattering is initiated only
by high-energy photons, like that of high-energy X-rays and gamma rays.
It was not possible to explain Compton scattering on the basis of classical picture of
electromagnetic radiations which assumes EM radiations as waves. Only the particle
nature or quantum aspect of EM radiations can explain the Compton scattering.
In quantum mechanical framework, a gamma ray is treated as quanta of energy
which behaves as a particle. Compton scattering, in quantum approach, may be
treated as the inelastic scattering of incident gamma ray photon with the free and
stationary electron in the target atom. Since the energy E (= hν) of the incident
gamma ray photon is much larger (≈ few hundred keV) than the binding energy
Be (≈ few eV) of the outer electrons of the target atom, it is reasonable to assume
that the outer shell electron of an atom is free (unbound) and stationary. Pictorial
representation of inelastic scattering of high-energy photon by the stationary and
free electron is shown in Fig. 4.29. As shown in the figure, a photon of wavelength
λi impinges on the stationary and free electron and kicks it out in a direction making
an angle φ with the incident direction with kinetic energy E kin
e
.
The incident photon gets scattered in a direction making angle θ with direction
of incidence, with reduced energy having a longer wavelength λf . One may apply
the laws of conservation of energy and linear momentum to the inelastic collision
between the photon and the stationary electron. However, one point must be kept in
mind while applying the law of conservation of energy it is that photo always move
with the velocity of light and that the ejected electron may be given high velocity
Fig. 4.29 Pictorial

representation of Compton
scattering
where relativistic variation of mass may become significant, therefore, one must use
m e c2 for the rest mass /
energy of electron and the energy of electron after collision

e
E kin may be given as m 2e c4 + pe2 c2 ; here m e is the rest mass of the electron, pe
the final linear momentum of electron and c the velocity of light.
Energy conservation
/
hνi + m e c2 = hνf + m 2e c4 + pe2 c2 (4.62)
Conservation of X-component of linear momentum
hνi hνf
= pe cos φ + cos θ (4.63)
c c
Conservation of the Y-component of linear momentum
hνf
0 = pe sin φ + sin θ (4.64)
c
Equations (4.62), (4.63) and (4.64) may be solved to get Compton equation.
h
λf − λi = Δλ = (1 − cos θ ) (4.65)
mec
Compton equation (4.65) tells that Compton shift in the wavelength Δλ can have a
minimum value of zero, when incident photon passes on along the incident direction
without getting scattered by electron, and the magnitude of wavelength shift increases
with the angle of scattering θ, attaining a maximum value (2h/me c) for backscattering
(θ = 180°) of photon.
In his original experiments, Compton bombarded carbon target with high-energy
X-ray photons and recorded the scattered photon of lower energy. Compton could
explain the experimental data assuming the particle nature of photon and inelastic
4.7 Compton Scattering of Gamma Rays 245
scattering of incident photon by stationary electron. It was the time when the particle
aspect of photon suggested by photoelectric effect was still being debated, Compton’s
analysis of his experiments gave a clear and independent evidence of particle-like
behaviour of electromagnetic radiations.
4.7.1 Compton Wavelength
The quantity mhe c is called the Compton wavelength of electron. In general the
Compton wavelength of a particle of rest mass m0 is given as mh0 c . Physical signif-
icance of Compton wavelength may be derived from the de Broglie wavelength
associated with a particle, according to which a particle of rest mass m0 moving with
velocity v has an associated de Broglie wavelength λde-bro = mh0 v . Since the velocity
of a moving particle cannot exceed the velocity of light c, the minimum value of
the de Broglie wavelength may occur when the velocity of the particle is taken as c.
Putting v = c, gives the de Broglie wavelength as mh0 c , which is the Compton wave-
length of the particle. Since wavelength associated with a particle is a measure of the
uncertainty in the position of the particle, a particle cannot be confined in a space
smaller or equal to its Compton wavelength. For example, Compton wavelength of
6.626×10−34 J s −12
electron = mhe c = 9.1×10 −31 kg×3.0×108 m = 2.427 × 10 m is around 2.4 × 10−12 m,
and therefore, an electron cannot be confined in a space equal or shorter than this.
Since the size (radius) of an average nucleus is of the order of 10−14 m, two order of
magnitudes is smaller than the Compton wavelength of electron: electron cannot be
confined within the nucleus and cannot be a constituent of nucleus.
4.7.2 Compton Scattering by the Whole Atom
In some experiments Compton recorded incident photons scattered by large angle

without appreciable change in wavelength. Also in such events no electron was
detected. Compton explained such events by assuming that the incident photon is
scattered not by the electron of the atom, rather it is scattered by the atom as a whole.
The change in wavelength Δλ in such scattering by atom may be written as:
h
λf − λi = Δλ = (1 − cos θ )
Mato c
M ato in the above expression is the mass of the atom as a whole, and Mhato c is
Compton wavelength of the atom. Since mass of the atom is very large, its Compton
wavelength is very small; hence, change in wavelength is undetectable. Further, the
target atom remains intact; no electron is ejected by the incident photon.
4.7.3 Photon Interactions with Matter
Photons, depending on their energy and the atomic number of the target atom, may
interact with the atom, with the bound and the free electrons of the atom and with
the nucleus of the target atom.
Low-energy photons mostly interact with the bound electrons (K-shell or L-shell
electrons) producing photoelectric effect. With the increase of energy the probability
of photon interaction with loosely bound electrons of the atom increases, resulting
in Compton scattering. In case the energy of the photon is larger than 1.02 MeV,
it may annihilate producing an electron and positron pair: the process called pair
production. The minimum energy of photon required to produce an electron positron
pair is 1.02 MeV which is the sum of the rest mass energies of the electron and positron
pair (0.51 + 0.51 = 1.02). Pair production takes place in the field of a nucleus,
when a high-energy photon (E phot > 1.02 MeV) passes through the nuclear field.
Nuclear field facilitates the recoil of the nucleus that is required for the conservation
of momentum in pair production process. High-energy photons may also excite or
disintegrate atomic nucleus. Very high-energy photons, with energies > 150 MeV,
may create mesons.
Figure 4.30 shows the variation with photon energy of the probability for photo-
electric effect, Compton scattering and pair production in lead (Pb). Kink (towards
the top) in the curve for photoelectric effect, called k-edge, shows that probability
for photoelectric effect suddenly increases for photon energy corresponding to the
binding energy of the K-shell electrons. Similar but less pronounced edges (not
shown in the figure) also appear for L and M-shells.
SAQ: What may be the order of magnitude for Compton wavelength of a neutron?
Fig. 4.30 Probability of

photoelectric effect,
Compton scattering and pair
production as a function of
photon energy in lead (pb)
4.8 Specific Heat of Solids 247
4.7.4 Some Applications of Compton Scattering
For the explanation and recording of the Compton effect, Compton was awarded a
share of the Nobel Prize in physics in the year 1927. Not only the Compton effect
represented the particle nature of light, but it is also important from the application
point of view. The Compton scattering is of importance in material science where
it is being used to get information regarding wavefunction of electrons in matter. It
is also of importance in the field of radiobiology and radiation therapy. Compton
scattering also has applications in X-ray astronomy and in getting signature of black
hole.
SAQ: In nature there are no free electrons, then why the Compton scattering is said
to take place with free electrons?
4.8 Specific Heat of Solids
Specific heat is defined as the amount of heat required to change the temperature of
unit mass of a substance by unit degree temperature. If ‘m’ kg of a substance is given
a heat energy of amount ΔQ which rises the temperature of the substance by Δθ,
then the specific heat of the substance is
1 ΔQ ΔQ
Specific heat = and the heat capacity =
m Δθ Δθ
Molar or atomic specific heat It is defined as the quantity of heat energy required
to raise the temperature of 1 kg mol or 1 kg atom of any substance by unit degree. It
is obvious that molecular or atomic specific heats are the products of the molecular
weight or the atomic weight with the specific heat of the substance. Molecular or
atomic specific heat of solids is generally denoted by C v . [In case of gases there
may be two types of molecular/atomic specific heats: at constant volume C v and at
constant pressure C p .] In the following discussion the term atomic specific heat will
be used which will also mean molar specific heat in case the solid is a compound
and not an element.
4.8.1 Dulong–Petit Law
French chemist Pierre Louis Dulong and French physicist Alexis Therese Petit in
1819 on the basis of their observation of atomic specific heat for large number of
solids gave an empirical law which states that ‘gram-atomic heat capacity (atomic
specific heat) of an element is a constant: that is, it is same for all solid elements, about
6 cal per g atom per °C and it is independent of temperature’. More than 60 elements
in solid state are found to have their atomic specific heat in the range of 5.38–6.93 cal/
g atom/°C with an average of 6.15. However, the gram-atomic specific heat for some
light elements like Silicon (gram-atomic specific heat 4.95) and diamond with gram-
atomic specific heat of 1.46 cal/g atom/°C does not follow Dulong–Petit law. Further,
it is found that the gram-atomic specific heat of solids depends on temperature and
approaches zero at absolute zero of temperature.
The MKS unit for atomic specific heat is J/kg atom/K and 1 J/kg atom/K = 4.2 ×
103 cal/kg atom/K. Therefore, the average value of atomic specific heat of 6.16 cal/
g atom/°C observed by Dulong–Petit is equal to 25.67 × 103 J/kg atom/K.
4.8.2 Obtaining Dulong–Petit Law on the Basis of Classical

Physics
Dulong–Petit law may be derived assuming that the classical law of equipartition
of energy of thermodynamics holds good. At absolute zero solids have a crystalline
structure in which atoms or molecules of the solid are at rest being held at their
place by mutual interaction. When energy in the form of heat is supplied to the
solid, the atoms or the molecules start vibrating around their mean position. If the
temperature is not very high, the vibratory motion has six degrees of freedom: three of
translatory motion (associated with kinetic energy) and three of vibrational motion
(associated with potential energy). Now, according to the law of equipartition of
energy, 21 kB T of energy is associated with each degree of freedom, and hence, the
total energy u associated with each atom (or molecule) at temperature T (K) is
u = 6 × 21 kB T = 3kB T (J). If AV denotes the Avogadro number, that is the number
of atoms/molecules in one kilo atom (or mole) of the solid, then the energy possessed
by 1 kilo atom of the substance is
U = AV u = 3kB AV T (4.66)
Here kB is Boltzmann constant (= 1.380 × 10−23 J/K) and Avogadro’s number
AV = 6.03 × 1026 atoms per kg
But kB AV = R, where R is the gas constant having the value R = 8.4 ×

103 J/kg atom/K.
Substituting kB AV = R in Eq. (4.66), one gets
U = 3RT (4.67)
And the atomic (or molar) specific heat CV is given as,
dU
CV = = 3R (4.68)
dT
4.9 Quantum Approach to Atomic Specific Heat of Solids 249
Equation (4.68) says that the atomic or molecular specific heat for all solids has
a fixed value of 3R (≈ 6.8 cal/g atom/°C) and is independent of temperature. This
confirms Dulong–Petit law.
4.8.3 Problems with Dulong–Petit Law
According to Dulong–Petit law (supported by classical physics) atomic specific

heat should be same for all solids and be independent of temperature. Experimental
measurements do not support both the above predictions. Experiments indicate that
atomic specific heat for metallic solids has a value near to 3R and also changes slowly
with temperature, but for non-metallic solids the magnitude of atomic specific heat
was quite away from 3R and that it also increases with the increase in temperature
approaching the Dulong–Petit value. However, the biggest challenge to Dulong–
Petit law comes from the variation of experimentally measured atomic specific heats
at lower temperatures (see Fig. 4.31); it was observed that atomic specific heat of
solids decreases rapidly with the decrease in temperature approaching a value zero
at absolute zero of temperature.
As such Dulong–Petit law, backed by classical theory, fails to explain the depen-
dence of atomic specific heat of solids on temperature and its sharp fall approaching
to zero at absolute zero.
4.9 Quantum Approach to Atomic Specific Heat of Solids
Initially Einstein in 1905 used the concept that a solid contains quantum harmonic
oscillators all having same energy to derive atomic specific heat of solids. Einstein’s
formulation correctly predicted the temperature dependence of atomic specific heat,
Fig. 4.31 Temperature

dependence of atomic
specific heat for some
elements in solid state
but quantitative agreement with experimental values was poor. Later, in 1912 Peter
Debye modified the concept of atomic oscillators all of same energy and included
oscillators with different values of quantized energies. Debye’s theory correctly
predicted both the temperature dependence and the magnitudes of atomic specific
heats of solids. Essentials of both Einstein and Debye theories are discussed in the
following.
4.9.1 Einstein’s Theory for Specific Heat of Solids
Einstein made following assumptions for absorption of heat energy by solids:

(a) Atoms in a solid at absolute zero are stationary at their equilibrium position
under the mutual interaction between different atoms. As such each atom at
absolute zero has zero or no energy.
(b) When heat energy is given to a solid, its atoms start vibrating with a characteristic
frequency ν, and characteristic frequency has a different value for each solid.
Each atom of a given solid vibrates with the same frequency ν.
(c) Each atom of a solid has three degrees of freedom, like the molecule of a perfect
gas.
(d) Each vibrating atom of the solid behaves like a Planck’s oscillator, its energy is
quantized, i.e. the energy of each vibrating atom is an integral multiple of hν.
Each atom has same value of energy which is an integer multiple of hν.
Under these assumptions, the mean energy per degree of freedom is not 21 kB T
as in classical law of equipartition of energy, instead the mean energy per degree of
freedom ∈ is given by:
hν
∈ = hν (4.69)
e kB T − 1
Above expression for average energy of an oscillator has been derived earlier (see
Eq. 4.60).
Since each atom has three degrees of freedom, energy associated with each atom
u becomes
3hν
u = hν (4.70)
e kB T − 1
Also, one kg atom of a solid contains AV number (Avogadro’s number) of atoms,

and energy U of 1 kg atom of solid is
3AV hν
U = AV u = hν (4.71)
e kB T − 1
And the atomic specific heat Cv is given by:

hν
dU hν 2 e kB T
Cv = = 3AV kB hν 2
dT kB T
e kB T − 1

2 hν
hν e kB T
= 3R hν 2 (4.72)
kB T
e kB T − 1

The quantity hν kB
has the dimensions of temperature and is called Einstein’s
temperature which is denoted by θE . Einstein’s temperature θE has a different value
for each solid.
Equation
(4.72) may be written in terms of Einstein’s temperature by
substituting kB = θE to get,
hν

2 θE
θE e T
Cv = 3R θ 2 (4.73)
T
eT −1
E
Expression given by Eq. (4.73) is called Einstein’s specific heat equation or

relation. Experimental values of atomic specific heat for a given solid at different
values of absolute temperature T are fitted in Eq. (4.73) to obtain the best fit value
of Einstein’s temperature θE from which the characteristic frequency ν for the solid
may be obtained; this value of the frequency is denoted by νE and is called Einstein’s
frequency of the solid.
4.9.2 Investigating the Temperature Dependence

of Einstein’s Equation
The temperature dependence of Einstein’s equation may be investigated in two limits:

when (a) θTE « 1 or (b) θTE ≫ 1.
θE
(a) At high temperatures T ≫ θE , θTE « 1; and term e T ≈ 1.
θE 2
And e T ≈ 1 + θTE + 2!1 θTE + · · ·
2
≈ 1 + θTE neglecting higher order terms θTE , and
θ
therefore, e T − 1 ≈ θTE .
E
θE
θ
Substituting e T ≈ 1 and e T − 1 ≈ θTE in Eq. (4.73), one gets,
E

2 θE 2
θE e T
θE 1
Cv = 3R θ 2 = 3R θE 2 = 3R
T T
eT −1
E
T
It may be observed that for temperatures T much higher than Einstein temperature
θE , atomic specific heat of solids approaches the Dulong–Petit value 3R.

(b) For the case of low temperatures when θTE ≫ 1

θE θE θE θE
e T −1 ≈e T as e ≫ 1 if
T ≫1
T
Equation (4.73) in this case reduces to,

2 θE 2 θE
θE e T
θE e T
Cv = 3R θ 2 ≈ 3R θ 2
T T
eT −1
E E
eT
2
θE 1
≈ 3R
θE

T
e T
Or
2
θE 1
Cv = 3R
θE
(4.74)
T
e T
2
In Eq. (4.74) on the right-hand side, there are two factors θTE and 1θE .
e T
θE 2
Factor T increases with the decrease of temperature T; however, the other
factor 1θE decreases exponentially with the decrease of temperature T. Rate of
T
e
decrease with temperature of the second factor is much faster as compared to the
rate of increase (with the decrease of temperature) of the first factor. As a result
the second factor becomes zero for very low temperatures earlier than the first factor
becomes infinite, hence, the atomic specific heat of solids approaches zero at absolute
temperature T approaches zero.
4.9.3 Drawbacks of Einstein’s Model
Though Einstein’s theory for specific heat of solids predicts that the atomic specific
heat for all solids should approach Dulong–Petit value of 3R at high temperatures
and it should approach zero at 0 K temperature, but it could not reproduce the
experimental values of specific heats for most of the solids. This theory suffers from
the following drawbacks:
(i) Does not reproduce the experimental values of specific heats for most solids.
(ii) Einstein temperature ϑE and Einstein frequency νE have no physical justifica-
tion; they could not be associated with any property of the solid, like its elastic
constants or melting point, etc.
4.9.4 Debye Theory of Atomic Specific Heat
Einstein in his theory for specific heat of solids assumed that on receiving heat energy,
each atom of the solid vibrates with the same frequency which is quantized. Debye,
on the other hand, assumed that on heating, the solid as a whole, i.e. the crystal
lattices in the solid, undergoes collective vibrations. Lattice vibrations are assumed
to be quantized. The quanta that represent lattice vibration are called PHONON.
Phonon is the counter part of photon which is the quanta of EM waves. In case of the
blackbody radiations it was assumed that the blackbody cavity is filled with photons
of different quantized frequencies; similarly, Debye assumed that on heating a solid
it gets filled with phonons of different frequencies that have quantized energies.
Often, it is said that a solid at some temperature above absolute zero is filled with
a phonon gas. In case of blackbody radiations it was assumed that photons of all
frequencies are present in the blackbody cavity. However, in case of vibrations in
a solids phonons of all frequencies are not present, and phonon frequency is bound
by the medium of its propagation which is the atomic lattice of the solid. So there
is an upper limit on the frequency of phonon in the solid that depends on the elastic
constants and the crystal structure of the material. Debye also assumed that phonon
waves, that are elastic waves, travel with some finite speed in the solid medium like
sound waves. On account of their interference, standing phonon waves are formed
in the solid. There may be three types of standing waves in the solid, longitudinal
waves of velocity C L and two types of transverse waves with two different states of
polarisations with speed C T .
Before proceeding further, let us point out the basic difference in Classical theory,
Einstein’s theory and Debye theory of specific heat of solids.
(i) Both the classical (Dulong–Petit) and Einstein’s models treat each atom of the
solid independently.
(ii) Debye model is more realistic, since it considers collective vibrations of many
atoms. As a matter of fact if one atom of the solid vibrates, the neighbouring
atoms are also set in vibratory motion.
(iii) Einstein’s model assumes that all atoms vibrate with the same frequency. Debye
model, on the other hand, considers the collective vibrations of different groups
of atoms. A group with large number of atoms can vibrate with lower frequency
while the group with fewer atoms may vibrate with higher frequency. Therefore,
the solid will contain vibrating groups of atoms with different frequencies, and
the frequencies are being quantized. Vibrating groups of atoms fill the solid
medium with elastic waves (sound waves) which is being confined within the
boundaries of the solid get reflected at boundaries and form standing waves. The
wavelength of an elastic wave depends on the number of atoms in the vibrating
group. The minimum wavelength λmin will correspond to the vibration of only
few atoms and will be related to the lattice constant of the crystalline structure.
It is obvious that the group of atoms that vibrates with minimum wavelength
will have largest value of the frequency of vibrations νmax . Further, if υ is the
speed of the elastic wave in the medium, then υ = λmin νmax . Here νmax is
the frequency of the wave that has the minimum wavelength. It follows from
here that for a given material there will be standing waves of several quantised
frequencies up to the maximum frequency νmax . Also, the νmax will depend
both on the crystal structure (lattice constants) and the speed of the elastic
waves in the medium. It is known that that elastic waves may be of two types:
longitudinal that may travel with some speed say C 1 and transverse that may
travel with a different speed, say C 2 in the medium. Further, transverse waves
may have two different states of polarisations; therefore, there may be three
different types of elastic waves of each frequency ν.
(iv) In Einstein model the energy E of the system is given by
E(Eins) = Average energy associated with each mode × number of atoms.
In Debye model system energy is given as

E(Debye) = Average energy associated with each mode of frequency ν
ν
× number of modes of frequency ν (4.75)
In Einstein model the number of modes is taken equal to the number of atom in
a kilomole (= Avogadro number AV ).
Modes essentially mean the number of standing waves in the volume of the solid.
Since the number of standing waves is quite large, one calculates the number of
standing waves g(ν)dν in a small frequency interval ν and (ν + dε) and integrates it
from zero to νmax to get the total numbers of modes.
It can be shown that the number of modes g(ν)dν for frequency range ν and
(ν + dν) for kg atom of the solid is given as:
9AV ν 2
g(ν)dν = (4.76)
νmax
3
Further, there are large number of phonons with different energies, and the average
energy Eave associated with each mode of vibration may be calculated using quantum
mechanical Maxwell Boltzmann statistics and is given as,
hν
Eave = hν
(4.77)
e kB T
−1
The system energy E ν corresponding to phonon of frequency ν may be calculated

by putting the values of different factors from Eqs. (4.76) and (4.77) in Eq. (4.75) to
get
⎛ ⎞

hν 9AV ν 2
E ν = (Eν )(g(ν)dν) = ⎝ hν ⎠ (4.78)
e kB T − 1 νmax
3
The total energy of the system may be obtained by integrating the above expression
over frequency ν
⎛ ⎞
∫νmax ∫νmax 2 ∫νmax
⎝ hν hν 9A ν 9A hν 3
⎠
V V
E= E ν dν = dν = 3 hν dν
e kB T
− 1 νmax
3 νmax e kB T
− 1
0 0 0
(4.79)
⎡ ⎤
⎢ 9A ∫ νmax ⎥
d⎢ ⎞ dν ⎥
3
V
⎣ νmax
3 0
⎛ hν
hν ⎦ hν −1
⎝e B −1⎠
k T
∫ νmax ∂ e kB T −1
But CV = dE
dT
= dT
= 9AV
νmax
3 0 hν 3 ∂T
dν
hν
−1
∫ νmax ∂ e kB T −1 ∫ νmax hν
hνe kB T
Or CV = 9AV
νmax
3 0 hν 3 ∂T
= 9AV
νmax
3 0 hν 3 hν 2 dν.
kB T e kB T −1
2
Or
⎡ ⎤
∫νmax hν
⎢ ν e
2
4 kB T
9AV h ⎥
CV = 3 ⎣ hν 2 ⎦dν (4.80)
νmax kB T 2
0 e kB T − 1
Let us make the following substitutions in Eq. (4.80):
hνmax hν h hνmax
= θD (Debye temperature) and x = ; dx = dν, xmax =
kB T kB T kB T kB T
Therefore,
∫xmax kB T 4 x
9AV h 2 x e kB T
CV = 3 h
dx
νmax kB T 2 (ex − 1)2 h
0
∫xmax ∫xmax
9AV T 3 kB4 x 4 ex 9AV kB T 3 x 4 ex
= 3 dx = 3 dx
νmax h 3 (ex − 1)2 hνmax (ex − 1)2
0 kB 0
Or
3 ∫xmax
T x 4 ex
CV = 9AV kB dx (4.81)
θD (ex − 1)2
0
Temperature dependence of the atomic (or molar) specific heat of solids in case
of Debye theory may be discussed through expression (4.81). The value of C V for
high temperatures T >> θ D and for the case of low temperatures T << θ D is discussed
here.
(a) At high temperatures T: x = hν
kB T
is small
2 3
Therefore ex = 1 + x + x2! + x3! + · · · ≈ 1 + x and (ex − 1) ≈ x and (ex ≈ 1).
Substituting (ex − 1) = x in Eq. (4.81), one gets:
3 ∫xmax 3 ∫xmax 4
T x 4 ex T x ·1
CV = 9AV kB dx ≈ 9AV kB dx
θD (e − 1)
x 2 θD (x)2
0 0
3 ∫xmax 3 3 xmax
T ! 2" T x
CV ≈ 9AV kB x dx = 9AV kB
θD θD 3 0
0
3 3
T θD
= 3AV kB
θD T
Or
CV ≈ 3AV kB ≈ 3R; (4.82)
since AV (Avogadro’s no.) × kB (Boltzmann constant) = R (gas constant).
Equation (4.82) shows that at high temperatures the molar or atomic specific heat
for all solids approaches the Dulong–Petit value.
∫ x 4 ex
(b) At low temperature, when θTD > 1; 0 max (exx −1) 2 dx ≈ 15 π .
4 2
And
3
T 4 2
CV ≈ 9AV kB π (4.83)
θD 15
It follows from Eq. (4.83) that at low temperatures, specific atomic or molar heat
decreases as the third power of the absolute temperature and approaches zero at
absolute zero.
It may be remarked that Debye theory correctly predicts the experimentally
observed behaviour of the atomic or molar specific heat of solids.
4.9.5 Debye Temperature θD
One big drawback of Einstein’s theory was that Einstein temperature θE was not
related to any property of the solid. Debye temperature θ D , on the other hand, depends
on νmax which in turn is related to the speed of the elastic wave in the solid medium.
Therefore, Debye temperature depends on elastic constants of the solid. Since the
minimum wavelength λmin or maximum frequency νmax may be correlated with the
lattice parameters of the crystal structure, the lattice constant of the crystal structure
may be derived from Debye temperature. For example the size of the vibrating units
sets a limit on the minimum wavelength since shorter wavelengths do not lead to
new modes. The smallest unit of a crystalline solid is the unit cell. Thus, the unit cell
puts constrain on the minimum wavelength of the vibration as:
a = λmin (4.84)
where ‘a’ is the length of the unit cell.

Further, Debye temperature of a solid is the temperature above which the system
may be described by classical properties. Debye temperature for some metals is:
Aluminium 433 K; Antimony 220 K; copper 343 K; Germanium 374 K; gold
170 K; iron 470 K; lead 105 K. It may be observed that soft metals like gold and
lead have a lower value of Debye temperature.
Solved Examples
SE4.1 What is the minimum wavelength of the produced X-ray if 8 kV potential

difference is applied across the anode and the cathode?
Solution:

hc 6.6 × 10−34 J s 3 × 108 m s−1
λminimum = =
eV 1.6 × 10−19 8 × 103 J
= 1.546 × 10−10 m = 1.55 Å
SE4.2 An X-ray beam of wavelength 0.075 nm gets diffracted by (111) plane of a

salt, having lattice constant of 0.32 nm. Calculate the glancing angle for the
second order diffraction.
Solution: According to Bragg’s law, we have
2d sin ϑ = nλ
It is given that λ = 0.075 nm = 7.5 × 10−11 m; order of diffraction n = 4.

The value of lattice constant, a = 0.32 nm = 3.2 × 10−10 m
a 3.2 × 10−10
d=√ =√ m
h2 + k2 + l 2 12 + 12 + 12
3.2 × 10−10
d= √ m
3
Now substituting the values in Bragg’s equation,

3.2 × 10−10
2 √ m sin ϑ = 2 7.5 × 10−11 m
3
√
3 × 2 7.5 × 10−11 m
sin ϑ = = 1.732 × 4.34 × 10−1 = 0.405
2 × 3.2 × 10−10
ϑ = sin−1 (0.405) = 23.89◦ ≈ 24◦
SE4.3 Photons of wavelength 400 nm are incident on a metallic surface resulting

into the emission of electrons of 1.2 eV kinetic energy. Calculate (i) the work
function Φ of the material and (ii) the longest wavelength of emitted photons
which may cause emission of electrons from the given metallic surface.
Solution:
(i) First, we need to calculate the energy of incident radiation as: E = hc
λ
.
h = 6.626 × 10−34 J s
c = 3.0 × 108 m/s
λ = 400 nm = 400 × 10−9 m
Substituting the values in the expression E = hc

λ
to get the value of energy
in Joules
6.626 × 10−34 J s × 3.0 × 108 m/s

E= = 4.97 × 10−19 J
400 × 10−9 m
Since we wish to convert the energy in eV, we use the relation 1 eV = 1.60 ×
10−19 J
4.97 × 10−19
E= eV
1.6 × 10−19
E = 3.10625 eV
Work function Φ = Energy of incident radiation – Kinetic energy of emitted

electron
Φ = 3.10625 − 1.2 eV
Φ = 1.90625 eV
(ii) The longest wavelength of emitted photons which may cause emission of elec-
trons from the given metallic surface will be if the kinetic energy of emitted
electron becomes almost zero. In that case one gets:
Work function Φ = Energy of incident radiation
i.e. Φ = hcλ
gives
Therefore, energy of incident radiation = 1.90625 eV = hc λ
.
Substituting the values of Plancks constant and speed of light, one gets the
wavelength of incident radiation to be nearly equal to 6.93 × 10−9 m.
SE4.4 Calculate the value of Compton wavelength for an electron and compare it
with the size of an average nucleus.
Solution: Compton wavelength = h
m0 c
,
where h = Plancks constant = 6.6 × 10−34 J s
m 0 = rest mass of electron = 9.1 × 10−31 kg
c = speed of light in vacuum = 3.0 × 108 m/s
Substituting these values, one gets the Compton wavelength equal to 242 ×
10−14 m.
On the other hand, the size of nucleus is typically ≈ 10−14 m.
SE4.5 Calculate the change in wavelength of 511 keV photons undergoing Compton
scattering at an angle of 60°.
Solution: Energy of photon (E) = 511 keV = hc/λ.
Wavelength of photon may be given by, λ = hc
E

6.6 × 10−34 J s 3.0 × 108 ms−1
λ=
511 keV

6.6 × 10−34 J s 3.0 × 108 m s−1
λ=
511 × 106 eV 1.6 × 10−19 J

6.6 × 10−34 J 3.0 × 108 m
λ=
817.6 × 10−13 J

19.8 × 10−13 m
λ=
(817.6)
λ = 0.024217 × 10−13 m
λ = 4.4217 × 10−15 m
6.6×10−34 J s
(i) Change in wavelength Δλ = − cos 60◦ )
(9.1×10−31 kg)·(3.0×108 m/s) (1
Δλ = 241.75 × 10−11 (1 − cos 60◦ )
Δλ = 241.75 × 10−11 (1 − 0.5)
Δλ = 241.75 × 10−11 (0.5)
Δλ = 120.87 × 10−11 m
SE4.6 What is the wavelength of the scattered photon in question SE4.5 above?
Solution: The change Δλ in the wavelength of scattered photon λ, − λ = 120.87 ×

10−11 m
λ, = 120.87 × 10−11 m + λ

λ, = 0.012087 × 10−15 + 4.4217 × 10−15 m
λ, = 4.4338 × 10−15 m
SE4.7 Calculate the amount of energy needed to heat the metallic ball of 450 g
from 30 to 80 °C. Given, the specific heat of metallic ball is 0.129 J/g°C.
Solution: Given,
m = 450 g.
c = 0.129 J/g°C.
ΔT = (Final Temperature − Initial Temperature)

= C(80 ◦ C − 30 ◦ C) = 50 ◦ C
Substituting these values in the equation of specific heat Q = mcΔT
Q = 450 g × 0.129 J/g◦ C × 50 ◦ C
Q = 2904.5 J
Thus, it requires 2904.5 J of energy to heat the metallic ball from 30 to 80 °C.
1. Show that the low-frequency limit of Planck’s Law reduces to the Rayleigh–
Jeans Law.
2. Obtain an expression for the energy of recoiling electron in case of Compton
scattering.
3. What is the maximum energy transferred to an electron in the Compton
scattering?
4. In case of photoelectric effect does the residual atom also recoil? If yes, what
is its effect on the energy of emitted photoelectron?
5. Plot a graph for the intensity of electrons emitted as a function of scattering
angle and explain its various features.
6. What is the angular distribution of electrons emitted in the photoelectric effect?
7. Show that the average forward momentum of electrons emitted in photoelectric
effect is larger than the momentum of incident photon.
8. For higher energy photons the emitted photoelectrons are mostly forward
peaked. Why?
9. State Moseley’s law and discuss its importance.
10. Which experiment conclusively proved wave nature of particles? Very briefly
describe the experiment.
11. Waves carry some disturbance or variation of some parameter with time. Which
parameter varies with time in case of matter waves?
12. What is the principle of working of an electron microscope? Why it has high
magnifying power and resolution?
13. What is the difference between phase and group velocities? Which of them
represents the actual motion of an associated particle?
14. What is (are) main point(s) of difference between Debye and Einstein theories
for atomic specific heat of solids?
15. What is the physical significance of Debye temperature?
16. There are generally two velocities associated with matter waves; what are these
velocities called and what do they represent?
17. Give at least one reason why electron cannot be a constituent of nucleus.
MC4.1 Wien’s displacement law fails at,

(a) Smaller wavelengths
(b) Larger wavelengths
(c) Low temperatures
(d) High temperatures
ANS: (a)
MC4.2 A distant star radiates with maximum emission in the wavelength region
from 630 to 800 nm. If the Wien’s constant is equal to 4.8 × 10−3 mK, then
corresponding star temperatures for maximum emission are approximately
equal to:
(a) 4444 K, 3500 K
(b) 3333 K, 3200 K
(c) 4555 K, 3656 K
(d) 5434 K, 3502 K
ANS: (a)
MC4.3 If the temperature of a blackbody increases from 400 to 600 K, the emission
rate of energy will change by:
(a) 2 times
(b) 4 times
(c) 8 times
(d) 16 times
ANS: (d)
MC4.4 For a given photosensitive surface the work function is 4.62 × 10−19 J. If
the value of Plancks constant is 6.602 × 10−34 J s, the threshold frequency
is approximately equal to:
(a) 6.99 × 1014 Hz
(b) 5.99 × 1014 Hz
(c) 3.79 × 1014 Hz
(d) 6.89 × 1014 Hz
ANS: (a)
MC4.5 Light of frequency ν, which is greater than threshold frequency ν0 is inci-
dent on a photosensitive material; the number of photoelectrons emitted
is proportional to:
(a) Intensity of light
(b) Frequency of light
(c) Threshold frequency
(d) Energy of incident radiation

ANS: (a)
MC4.6 Stopping potential for photoelectric emission from a surface is 20 V: the
maximum kinetic energy of emitted electrons is nearly equal to:
(a) 4.3 × 10−18 J
(b) 6.2 × 10−18 J
(c) 8.2 × 10−18 J
(d) 3.2 × 10−18 J
ANS: (d)
MC4.7 Which of the following statement holds good for a blackbody?
(a) A perfect absorber and a perfect radiator
(b) A perfect absorber but not a perfect radiator
(c) A perfect radiator but not a perfect absorber
(d) A perfect insulator and a perfect conductor
ANS: (a)
MC4.8 The de Broglie wavelength of an electron accelerated from rest to a
potential of 500 V is,
(a) 0.5250 Å
(b) 0.3545 Å
(c) 0.5486 Å
(d) 0.7569 Å
ANS: (c)
MC4.9 X-rays having wavelength 10 Å are scattered from an Aluminium cylinder.
If the scattered radiations are observed at right angles to the direction of
incident X-ray beam, the Compton shift is approximately equal to?
(a) 0.0019 nm
(b) 0.0039 nm
(c) 0.0024 nm
(d) 0.0048 nm
ANS: (c)
MC4.10 In an experiment of Compton scattering the photon of energy hν is incident
on a free electron such that
the recoiling
electron has an energy given by
G(1−cos ϑ)
the expression, Re = hν 1+G(1−cos ϑ) . The ter constant G is equal to:
hν
(a) m0
hν
(b) m0c
hν
(c) m 0 c2
m 0 c2
(d) hν
ANS: (c)
LA4.1 Explain, why wave nature of light could not explain the phenomenon of
photoelectric effect.
LA4.2 Differentiate between characteristic and continuous X-rays. Explain on
what factors the intensity of Bremsstrahlung radiations depend.
LA4.3 Explain how one can differentiate between electrons emitted as a result of
Auger and beta decay process.
LA4.4 Calculate the frequency (in Hz) of X-ray emitted, when an atom de-excites
from a level of energy 662 to 300 eV.
LA4.5 Calculate the de Broglie wavelength of electron of energy 500 keV.
LA4.6 Show that the Rayleigh–Jeans law and the Wein’s law are the special cases
of Einstein’s law for blackbody radiation formula.
LA4.7 Calculate the energy of scattered photon undergoing Compton scattering
at 60°.
LA4.8 A body at 47 °C radiates photons. If the Wein’s constant is 4.898×10−3 mK,
what will be the peak of the wavelength radiated?
LA4.9 Show that at very large wavelengths λ, the Planck’s formula for spectral
radiations,

8π hc 1
E(λ) = hc
λ5 exp λkT −1
Reduces to the Rayleigh–Jeans law, given as:
8π
E(λ) = kT.
λ4
LA4.10 What is meant by the dual nature of matter? With necessary details describe
an experiment that conclusively proved the wave nature associated with a
particle.
LA4.11 What are X-rays, how are they produced? What is Moseley’s law? Give
properties of X-rays and clearly distinguish between continuous and
characteristic X-rays.
LA4.12 What are gamma rays? In what respect they are different from X-
rays? Under what conditions an electromagnetic radiation will produce
photoelectric effect and Compton scattering.
LA4.13 Discuss the phase and group velocities, how do they differ from each other?
LA4.14 Briefly discuss three failures of classical physics.

LA4.15 Discuss in details Einstein and Debye theories for the molar specific heat
of solids. What are the main points of difference between the two theories?
Chapter 5
Introduction to Quantum Mechanics
Objective
An introduction to Schrodinger picture of quantum mechanics is presented in this
chapter. Postulates of quantum mechanics, definitions of operators, operator algebra,
hermitian operators, eigen values and normalised eigen states, etc., are explained in
simple language and with some examples. Application of quantum mechanics has
been explained taking some examples of one-dimensional potentials. It is expected
that after reading this chapter a reader will be able to apply quantum mechanics to
simple cases.
5.1 Introduction
Physics attempts to describe the objects, systems and events of inanimate world and
their time evolution using mathematical language of functions and equations. Further,
the inanimate world is considered to have two components: the matter and the radi-
ations. Matter was assumed to have been made up of particles, while radiations were
considered to have wave nature. Till early twentieth century, behaviour and time
evolution of processes associated with matter were treated using Newton’s laws of
motion and those of radiations using Maxwell’s equations. Newtonian mechanics and
Maxwell equations are the two pillars of classical physics. Classical physics, to a
large extent, successfully explained most of the phenomena of inanimate world partic-
ularly at macroscopic level. However, problems become apparent with the advent of
experimental tools that opened the microscopic world for investigation. It was soon
realised that matter-radiation approach of classical physics is not adequate to explain
phenomenon at microscopic level. The classical theory failed to explain discrete
energy levels of atomic states, why atom is spherical and the wave particle duality,
energy distribution in blackbody radiations, photoelectric effect, Compton scattering,
specific heat of solids, etc., as discussed in Chap. 4. Failure of classical physics,
particularly in explaining physics at the microscopic level, accentuated the necessity
https://doi.org/10.1007/978-3-031-32084-2_5
268 5 Introduction to Quantum Mechanics
of some alternate physics which may satisfactorily explain physics of microscopic

world and may converge to classical physics in case of macroscopic world. Many
random and unrelated attempts were made, but the real breakthrough came in 1925
when Werner Heisenberg developed the first version of quantum-matrix mechanics
using which it became possible to calculate the position and intensity of spectral
lines emitted from hot glowing hydrogen gas. Heisenberg’s matrix representation of
quantum mechanics was difficult to follow and provided no insight into the physical
process underlying the process of discrete energy lines.
In 1926 Erwin Schrodinger published what is now known as the Schrodinger equa-
tion for quantum mechanics. Schrodinger was trying to generalise one-dimensional
approach of De Broglie for hydrogen atom (that could not explain why atom should
be spherical) to a three-dimensional picture, in which a 3-dimensional standing wave
spreads spherically around the core. He succeeded in finding the equation of such
a wave, which is called Schrodinger wave equation. Heisenberg and Wolfgang
Pauli in 1929 published the foundations of relativistic quantum mechanics, with
which the theoretical basis of quantum mechanics was put on a firm footing. Later,
it was shown that both the matrix approach of Heisenberg and the wave approach
of Schrodinger lead to the same results and are alternate descriptions of the same
physics. The main distinguishing features of Schrodinger and Heisenberg pictures
of quantum mechanism may be summarised as follows.
In Schrodinger picture, (i) quantum systems are regarded as wavefunctions that
are solutions of the Schrodinger equation, and (ii) observables are represented by
hermitian operators that act on wavefunction. (iii) In this picture operators stay fixed,
while Schrodinger equation changes the basis with time.
In Heisenberg picture, (i) operators change in time while the basis of space remains
fixed. A fixed basis, in some ways, is mathematically more elegant, and this formalism
easily generalises to relativistic case.
There is a third picture of quantum mechanics also put forward by Paul Dirac, in
which both the operators and the basis carry time dependence. This picture forms
the basis of quantum field theory.
In this chapter we will discuss only the fundamentals of Schrodinger’s picture of
quantum mechanics. Further, the term system will be used to represent a particle or
a group of particles or any other assembly of entities, throughout this chapter.
5.2 Postulates of Quantum Mechanics
In classical mechanics one defines the initial state of a particle or a system in terms
of its position and linear momentum at some initial time. These two parameters
completely define the state of the particle/system in classical mechanics. Next one
uses Newtonian mechanics to study the time evolution of the particle/system, i.e.
how different observable parameters, like the speed, linear momentum, energy, etc.,
of the particle/system change with time. Similarly, in quantum mechanics the first
5.2 Postulates of Quantum Mechanics 269
thing to do is to define the state of a system at some time say t. The first postulate of
quantum mechanics tells how to define the state of a system.
Postulate 1 The state of a quantum system at a given instant of time, t, is completely
defined by a function, ψ(→r , t), of position →r = {x, y, z} and time t.
Function ψ(→r , t) is called the wavefunction or the state function of the system
and is a complex-valued function that contains information about the position of the
system at time t.
5.2.1 What Does Wavefunction Represent?
No physical meaning may be assigned to the wavefunction ψ(→ r , t). However, a

statistical interpretation may be given. Suppose if one tries to find the system at place
−
→r , at time t in an infinitesimal small volume dV (= dx·dy·dz), then the probability
p of finding the system in this volume will be given by
( ⇀ )∗ ( ⇀ ) | ( ⇀ )|2
| |
p(r, dV, t) = ψ r , t ψ r , t dV = |ψ r , t | dV (5.1)
It follows from Eq. (5.1) that the probability P of finding the system in volume V
will be given by
˚ ˚ (⇀ )
P(r, t) = p(r, dV , t)dV = ψ r , t |2 dV (5.2)
If the integral on right-hand side of Eq. (5.2) is taken over the whole universe
that is over all the available space, naturally the probability of finding the system
somewhere in the universe will be 1. Therefore,
˚
+∞
(⇀ )
P(some where in the universe) = ψ r , t |2 dV = 1 (5.3)
−∞
Equation (5.3) is known as the normalisation( of the

) wavefunction. It may be shown
⇀
that Eq. (5.3) will hold only when function ψ r , t vanishes at r = ±∞.
Note: Though a volume integral is represented by three integration signs as has been
done so far but now on volume integrals will also be denoted by a single integration
sign to save space.
5.2.2 Properties of the Acceptable Wavefunction
Any function of r and t cannot be a valid wavefunction that represents the quantum
system at time t. A valid wavefunction must have the following properties.
(⇀ )
(i) The function must be single valued. A valid ψ r , t must have only one
value for the given values of r and t.
(ii) It must be continuous in the entire region defined by its independent variables.
(iii) Wavefunction must be finite everywhere, and the function and its derivatives
should vanish at infinity.
˝ (⇀ ) 2
(iv) It must be square integrable. Since ψ r , t | dV gives the probability of
finding the system at point defined by r at time t, and probability may have
some value between zero and 1, the square of the function must be integrable.
A function of (r, t) which does not satisfy even one of the above conditions cannot
be a valid quantum wavefunction of a quantum system.
5.3 Observables and Operators
There are always some measurable quantities that are associated with any system
and either all or some of these measurable quantities may change with time as the
system evolves. Let us take the example of a classical system of a particle of mass M
moving with velocity V which changes with time. In this case the linear momentum,
the angular momentum and the velocity are all measurable quantities that change with
time; therefore, they are dynamical measurable variables associated with the system.
In case of a quantum mechanical system the observables are treated as described by
the second postulate of quantum mechanics given below.
Postulate 2 An observable, A, is represented by a linear and hermitian operator
written as Â. Any operator, say, Â, is a mathematical instruction which when
applied to a mathematical object like the wavefunction ψ(x 1 , x 2 , x 3 . . .) gives
another similar mathematical object of same nature φ(x 1 , x 2 , x 3 . . .). It may be
noted that the object φ depends on same variables as the initial object ψ ; i.e.
both φ and ψ are functions of the same space. In mathematical language it may
be expressed as
Âψ(x 1 , x 2 , x 3 . . .) = φ(x 1 , x 2 , x 3 . . .) (5.4)
As mentioned, an operator is an instruction to carry out a particular mathematical Δ
operation on the given function. For example, an operator represented by M may

be the instruction to multiply an object χ (y1 , y2 ) by a constant K to get the object
ϕ(y1 , y2 ). This may be written as
5.4 Time Evolution of a Quantum Mechanical System 271
Δ
M χ (y1 , y2 ) = ϕ(y1 , y2 ) = K χ (y1 , y2 ) (5.5)

Δ
Equation (5.5) defines the eigenvalue K of operator M . In words one may say, if an
Δ
operator M operating on a wavefunction multiplies the wavefunction by a constant

Δ
K, the constant K is called the eigenvalue of operator M for the given wavefunction.
We shall discuss operators in more details in the following sections.
SAQ: What are the distinguishing characters of three versions of quantum
mechanics put forward, respectively, by Schrodinger, Heisenberg and Dirac?
5.4 Time Evolution of a Quantum Mechanical System
Time evolution, i.e. how the system changes with time, of a classical system is
governed by Newton’s laws of motion and Maxwell’s equations of electromagnetic
waves. The time evolution of a quantum mechanical system may be described by a
partial differential equation called Schrodinger equation.
Postulate 3 Time evolution of the wavefunction ψ(→r , t) of a quantum mechan-
ical system is governed by the following partial differential equation, called
Schrodinger time-dependent equation.
5.4.1 Schrodinger Time-Dependent Equation
∂ψ(→r , t) ℏ2 −
→2 (→ )
iℏ =− ∇ ψ(→r , t) + V −
r , t ψ(→r , t) (5.6)
∂t 2m
−r , t) the
In Eq. (5.6) m is the mass of the quantum mechanical system, V (→
potential energy function and,
−
→2 ∂2 ∂2 ∂2
∇ = + + , the Laplace operator.
∂ x2 ∂ y2 ∂ z2
In case of one spatial dimension, Schrodinger time-dependent equation reduces

to
∂ψ(x, t) ℏ2 ∂ 2 ψ(x, t)
iℏ =− + V (x)ψ(x, t) (5.7)
∂t 2m ∂ x 2
Solution of Schrodinger equation provides the wavefunction of the quantum
mechanical state.
5.4.2 Some Properties of Schrodinger Equation
(i) Schrodinger equation contains terms with only the first power of ψ; hence, it is
linear in ψ. Further, if a given function ψ is a solution of Schrodinger equation
for a given system, then βψ, where β is constant, will also be a solution of
the Schrodinger equation. This property is referred as homogeneity of the
equation. Schrodinger equation is, therefore, both linear and homogeneous.
(ii) Since Schrodinger equation is first-order differential equation in time, the value
of wavefunction ψ if known at some initial time t 0 , it may be determined at
some later time t uniquely.
(iii) Let us consider a collection of N number of particles in a small volume dτ
of space at some point r→ in space which we take as our quantum mechanical
system. Further, let us assume that each of these N particles is represented by
the same wavefunction ψ(→ r , t). The average number of particles at time t in
volume dτ , which is the probability of finding the collection of particles in
volume dτ , is given by
r , t)∗ ψ(→
⟨Nt ⟩ = N ψ(→ r , t)dτ (5.8)
Equation (5.8) may be used to calculate the average number of particles at different
points in the space at the same instant t by changing the value of position vector r→.
Thus, one may generate a probability distribution function of the system of particles
in space at time t. Same technique may be used to generate the probability distribution
function at some time t , and so on. In this way, Schrodinger equation may be used
to follow the time evolution of the quantum mechanical system of particles.
The significance of the time-dependent Schrodinger equation lies in the fact
that it allows the determination of the time evolution of wavefunction ψ which in
turn provides the time evolution of the probability density distribution function.
Changes in probability density distribution function tell about the changes that
have taken place in the system as it evolves with time.
(iv) If ψ1 (→
r , t), ψ2 (→
r , t) and ψ3 (→
r , t) are three solutions of the Schrodinger equa-
tion for a given system, then ψ(→ r , t) = a1 ψ1 (→
r , t) + a2 ψ2 (→
r , t) + a3 ψ3 (→
r , t),
where a1 , a2 and a3 are arbitrary complex constants, will also be a solu-
tion of the Schrodinger equation. This property of the Schrodinger equation
is called the property of quantum mechanical superposition. In general, the
quantum mechanical principle of superposition may be stated as: if there are
more than one wavefunctions that are solutions of the Schrodinger equation for
a given system, then the linear combination of these wavefunctions will also be
a solution of the Schrodinger equation.
SAQ: What is the counterpart of linear momentum in quantum mechanics?

5.5 Time-Independent Schrodinger Equation 273
5.5 Time-Independent Schrodinger Equation
Potential V in Schrodinger equation might be time dependent V (→ r , t) or it may be

time independent V (→ r ). In the case when potential does not depend on time, it is
possible to separate Schrodinger equation into two parts; one carrying terms having
time dependence and the other that does not depend on time, as shown here.
In case V is time independent, it is possible to write the state function ψ(→ r , t) as
ψ(→
r , t) = φ(→
r ) f (t) (5.9)
Substituting above expression for wavefunction ψ in Schrodinger equation gives
∂φ(→
r ) f (t) ℏ2 −
→2
iℏ =− ∇ φ(→
r ) f (t) + V (→
r )φ(→
r ) f (t)
∂t 2m
or
d f (t) ℏ2 −
→
φ(→
r )iℏ =− f (t) ∇ 2 φ(→
r ) + V (→
r )φ(→
r ) f (t) (5.10)
dt 2m
Dividing both sides of Eq. (5.10) by φ(→

r ) f (t) one gets
1 d f (t) ℏ2 1 − →2
iℏ =− ∇ φ(→
r ) + V (→
r) (5.11)
f (t) dt 2m φ(→
r)
It may be marked that the terms appearing on the RHS of Eq. (5.11) are all
functions only of r and term on the LHS is a function only of time t. It is obvious
that a time-dependent term cannot be equal to the sum of two time-independent
terms, as is seen in Eq. (5.11). This is possible only if the time-dependent and the
time-independent terms are equal to some constant, say E, i.e.
1 d f (t)
iℏ =E (5.12)
f (t) dt
and
ℏ2 1 − →2
− ∇ φ(→
r ) + V (→
r) = E (5.13)
2m φ(→
r)
Equation (5.12) may be written as
d f (t) i
= − E f (t)
dt ℏ
This equation may be easily integrated to give
f (t) = e(− ℏ Et )
i
(5.14)
Equation (5.13) may also be written as,
ℏ2 −
→2
− ∇ φ(→
r ) + V (→ r ) = Eφ(→
r )φ(→ r) (5.15)
2m
Equation (5.15) is called Schrodinger’s time-independent equation which may
be solved to obtain function φ(→
r ) provided the potential V (→
r ) is known.
Schrodinger’s time-independent equation may also be written as
Ĥ φ(→
r ) = Eφ(→
r) (5.16)
where
ℏ2 −
→2
Ĥ = − ∇ + V (→
r) (5.17)
2m
It may be noted that Ĥ in above equation is a differential operator which is

ℏ2 → 2
equivalent to the classical Hamiltonian H = − 2m ∇ + V (→ r ).
It follows from the derivation given above that in case the potential is time
independent, the total wavefunction ψ(→ r , t) may be written as
r )e(− ℏ Et )
i
ψ(→
r , t) = φ(→
r ) f (t) = φ(→ (5.18)
The probability density p for such case (where potential does not depend on time)
may be calculated using the expression
| | || | |
| |
|
|
p = |ψ(→ r )e(+ ℏ Et ) φ(→
r , t)| = |φ(→
r , t)∗ ψ(→ r )e(− ℏ Et ) | = |φ(→
i i
r )2 | (5.19)
Equation (5.19) tells that for a time-independent Schrodinger equation the prob-
ability density at a given point in space does not depend on time; it remains same at
all times.
The states of a quantum system for which the probability density does not depend
on time are called stationary states.
SAQ: Discuss the conditions when Schrodinger’s time-independent equation may
be used.
5.6 About Operators
Operators are instructions of mathematical nature that are required to be carried out
on the object (wavefunction) over which they are made to operate. Some important
classes of operators are discussed in the following.
5.6 About Operators 275
5.6.1 Null Operator (O)
A null operator, denoted by O, when operates on an object gives a zero; i.e. O

ψ(→
r , t) = 0.
5.6.2 Unity or Identity Operator ( Î)
The unity operator operating on an object leaves the object unchanged; for example,
r , t) = φ(→
Îφ(→ r , t).
5.6.3 Linear Operator
If, for{ the given scalars α and} β and functions ψ(→r , t) and φ(→ r , t),
A αψ(→ r , t) + β̀φ(→
r , t) = α Aψ(→r , t) + β Aφ(→r , t), then operator A is called a
linear operator.
5.6.4 Hermitian Conjugate and Hermitian Operator
Let there be a hermitian operator Â, we define another operator Â† such that
{+∞ [ ] {+∞( )∗
∗
φ (→ r )d x =
r ) Âψ(→ 3
Â† φ(→
r ) ψ(→
r )d3 x (5.20)
−∞ −∞
Then operator Â† is called the hermitian adjoint or hermitian conjugate of operator
Â.
Suppose the operator Â may be represented by a matrix M of appropriate dimen-
sions, then the hermitian conjugate operator will be given by the matrix M † which
may be obtained by first transposing the matrix M and then taking the complex
conjugate, i.e.
[ ]∗
M † = (M)T (5.21)
Operators that are their own hermitian conjugate are called hermitian operators;
that is to say that if Â† = Â, then Â is a hermitian operator.
One important property of hermitian operators is that their eigen values are
real. Since dynamic observables of quantum mechanical systems must also be real,
therefore, operators that represent dynamic variables must be hermitian.
5.6.5 Anti-hermitian Operator
An anti-hermitian operator is equal to the negative of its hermitian conjugate. Anti-

hermitian operator of hermitian operator Â is often denoted by (− Â) Therefore,
( ) †
− Â = − Â (5.22)
5.6.6 Inverse Operator ( Â−1 )
Inverse of an operator Â−1 is defined by the relation
Â−1 Â = Â Â−1 = Iˆ (unit operator) (5.23)
5.6.7 Unitary Operator (Û)
A linear operator is unitary if it satisfies the relation.
Iˆ
Û Û † = Û † Û = Iˆ or Û † = = Û −1 (inverse operator) (5.24)
Û
Equation (5.24) tells that the inverse of a unitary operator is equal to its hermitian
conjugate.
SAQ: What is meant by a linear operator?
5.6.8 Some Properties of Hermitian Operators
Without any derivation or proof we list here some properties of hermitian operators.
(i) The eigenvalues of a hermitian operator are real.
(ii) The eigenfunctions of a hermitian operator, corresponding to distinct eigen-

values, are orthogonal.
(iii) The eigenvalues of an anti-hermitian operator are either purely imaginary or
equal to zero.
(iv) A matrix operator Â has an inverse only if the matrix representing the operator
is a square matrix and its determinant is nonzero.
(v) The eigenvalues of a unitary operator are complex numbers of moduli equal
to one and the eigenfunctions of a unitary operator that does not have any
degenerate eigenvalue are mutually orthogonal.
5.6.9 Algebra of Operators
Definitions of different algebraic operations carried over operators will be explained

in this section.
(a) Sum of two operators: Suppose a hermitian operator Ĉ is given as the sum of
two hermitian operators Â and B̂, i.e. Ĉ = Â + B̂, then
( )
Ĉψ(x) = Â + B̂ ψ(x) = Âψ(x) + B̂ψ(x) (5.25)
(b) Products of an operator with a complex number: The product of an operator

Â with some complex number c, i.e. c Â is defined as
Δ
c Aψ(x) = c Âψ(x) (5.26)
(c) Product of two operators: An operator Ĉ may be defined as the product of two
operators Â and B̂ when Ĉψ(x) = Â B̂ψ(x) = φ(x).
Also,
B̂ψ(x) = ϑ(x) and Â B̂ = Âϑ(x) = φ(x) (5.27)
As is evident, operator B̂ operates on ψ(x) to give a new function ϑ(x) such that
when Â operates on ϑ(x) it gives the function φ(x).
(d) Division of two operators: Dividing an operator Â by another operator B̂ is
equivalent to multiplying operator Â by operator B̂ −1 , the inverse operator of
B̂, provided the inverse operator of B̂ exists.
(e) Power of an operator: Sometimes one finds expression like B̂ n ψ(x) where an
operator is raised to some power n. It simply means that the operator should
operate on the function successively n-times one after the other;
(n−1)
. .. .( )
B̂ n ψ(x) = B̂ B̂ B̂ . . . . . . B̂ B̂ψ(x)
(f) Commutator operator: In general the products of the two operators Â B̂ and
B̂ Â are not equal, i.e. Â B̂ /= B̂ Â. An operator Ĉ defined by the expression
(5.28) and written as [ Â, B̂] is called the commutator of Â and B̂.
[ ]
Ĉ = Â, B̂ = Â B̂ − B̂ Â (Commutator) (5.28)
In case the commutator of Â and B̂ is zero, i.e. Ĉ = [ Â, B̂] = Â B̂ − B̂ Â = 0, it

is said that the two operators Â and B̂ commute with each other.
Similarly, in general Â B̂ + B̂ Â /= 2 Â B̂ and in such cases one may also define an
anticommutator operator denoted as [ Â, B̂]+
[ ]
Â, B̂ = Â B̂ + B̂ Â (Anti-commutator) (5.29)
+
When anticommutator operator [ Â, B̂]+ = 0, the operators Â and B̂ are said to
anticommute with each other.
When a system defined by a wavefunction, say, ψ(x) is operated by the product
of two operators Â and B̂ and the final result of the operation does not depend on the
order of operations, the operators Â and B̂ are said to commute with each other. Just
for example, if operator Â means the instruction of writing your name and operator B̂
means drawing of an apple, then the order of operators is not important, the end result
of writing name and drawing the figure of an apple may be achieved by first writing
the name and then drawing the figure of apple, or first drawing the figure of apple and
then writing the name. In this example operator Â and operator B̂ commute. However,
if operator Â is the command to colour the apple and operator B̂ is the command to
draw the figure of an apple, then the product Â B̂ will be meaningless, while product
B̂ Â will mean to draw the figure of an apple and to colour it, a meaningful command.
In case operators Â and B̂ represent some dynamic variables of the system (like
the position and linear momentum, etc.) and if they commute, it means that both
dynamic variables of the system may be measured simultaneously with full accuracy
in the given state of the system. On the other hand if they do not commute, then both
the variables cannot be measured with complete accuracy in the given state of the
system. If two operators commute, then they can have the same set of eigenfunctions.
(g) Some properties of commutator: It may be proved that commutator obeys
following rules:
[ ] [ ]
(i) Â, B̂ = − B̂, Â
[ ]† [ † † ]
(ii) Â, B̂ = Â , B̂
[ ( )] [ ] [ ]
(iii) Â, B̂ + Ĉ = Â, B̂ + Â, Ĉ
[ ( )] [ ] [ ]
(iv) Â, B̂ Ĉ = B̂ Â, Ĉ + Â, B̂ Ĉ.
5.6.10 Operators for Some Dynamical Variables
Table 5.1 gives the list of some important and frequently used quantum mechanical
operators corresponding to classical observables.
SAQ: Why hermitian operators are used to describe observables?
Solved Examples
SE5.1 Show that function ϕ(y, t) = Aye−β y e−iωt where A, β and ω are arbitrary
2
constants, is an acceptable wavefunction for a quantum mechanical system.

Solution: Function ϕ(y, t) is single valued and continuous over all values of y and
t. The next point to check is if the function is square integrable or not. So, let us
evaluate the integral
{+∞ {+∞( )( )
∗
A∗ ye−βy e+iωt Aye−β y e−iωt dydt
2 2
I = ϕ (y, t)ϕ(y, t)dydt =
−∞ −∞
{+∞
∗
y 2 e−2β y dy
2
I =A A (5.30)
−∞
Table 5.1 Quantum

Operators in quantum Corresponding dynamical
mechanical operators for
mechanics variables in classical
some classical variables
mechanics
Radial coordinate r̂ r→
Cartesian coordinates x, y, z X, y, z
Δ Δ
Angular momentum r→ × p→ L→
( ) ( )
∂
L x = −iℏ y ∂z − z ∂∂y L x = ypz − zp y
( )
L y = −iℏ z ∂∂x − x ∂∂z L y = (zpx − x pz )
( ) ( )
L z = −iℏ x ∂∂y − y ∂∂x L z = x p y − ypx
p2
Total energy: 2m + V (→
r)
ℏ →22
H = − 2m ∇ + V̂ (→
r)
Parity operator Mirror reflection
−→
P = −r (−x, −y, −z)
P→ = −iℏ∇;→ px = −iℏ ∂
∂x ; Linear momentum P
py = −iℏ ∂∂y ; pz = ∂
−iℏ ∂z
ℏ →2 2 p2
Kinetic energy T = − 2m ∇ 2m
To evaluate integral I given by Eq. (5.30), we make use of the Gaussian integral
IGauu given by
{+∞ /
−αx 2 π
IGauu (α) = e dx = (5.31)
α
−∞
We differentiate Eq. (5.31) with respect to α to get
{+∞ /
dIGauu (α) d −αx 2 d π
= e dx =
dα dα dα α
−∞
or
{+∞ /
dIGauu (α) 2 −αx 2 π
= x e dx =
dα 4α 3
−∞
∫ +∞ √ π
x 2 e−αx dx =
2
or . Putting x = y and α = 2β in this equation one gets
∫−∞
+∞
/4α3 /
2 −2βy 2 π π
−∞ y e dy = 4(2β)3 = 32β 3 , and putting this back in Eq. (5.30) one get,
{+∞ /
∗ π
y 2 e−2β y dy = A∗ A
2
I =A A (5.32)
32β 3
−∞
Equation (5.32) shows that the function ϕ(y, t) is square integrable and hence
the given function fulfils all the conditions that a valid wavefunction of a quantum
mechanical system must fulfil. Therefore, it may be the state function of a quantum
system.
[ ]
SE5.2 Calculate the value of the commutator Â, B̂ when operators are defined
by the relations; Âϕ(x) = xϕ(x) and B̂ϕ(x) = −iℏ dϕ(x)
dx
, here ϕ(x) is a
valid wavefunction.
[ ] { }
Solution: The commutator Â, B̂ ϑ(x) = Â B̂ϑ(x) − B̂ Âϑ(x) = Â −iℏ dϕ(x)
dx
−
B̂{xϕ(x)}.
Or
[ ] ⎧ ⎫
dϕ(x) dxϕ(x)
Â, B̂ ϑ(x) = −iℏ Â − −iℏ
dx dx
( )
dϕ(x) dxϕ(x) dϕ(x) dxϕ(x)
= −iℏx + iℏ = −iℏ x −
dx dx dx dx
( )
dϕ(x) dϑ(x)
= −iℏ x −x − ϑ(x) = iℏϑ(x) (5.33)
dx dx
[ ]
Equation (5.33) tells that the commutator Â, B̂ = iℏ.
Note: It is easy to identify that operator Â corresponds to the x-component of the

position operator r̂ and operator B̂ to the x-component of linear momentum operator
Δ
px . The result given by Eq. (5.33) may be generalised to the relation

[ ]
x̂ j , p̂k = iℏδ j,k here δ j,k = 1 for j = k and δ j,k = 0 if j /= k (5.34)
δ j,k is called Kronecker delta function.
SE5.3 How one defines the hermitian conjugate of a hermitian operator? An oper-
ator given by Â = − d dϕ(y)
2
2
y
operates on a well-behaved state function ϕ(y).
Obtain an expression for the operator Â† and show it is same as Â.
Solution: Hermitian conjugate or adjoint of a hermitian operator is defined through

the following integral relation given by Eq. (5.20);
{+∞ [ ] {+∞( )∗
∗
φ (→ r )d x =
r ) Âψ(→ 3
Â† φ(→
r ) ψ(→
r )d3 x
−∞ −∞
In the present case, the operator acts on a wavefunction that is the function of
only y, the above expression may be written as
{+∞ [ ] {+∞( )∗
∗
φ (→y ) Âψ(→y )dy = Â† φ(→y ) ψ(→y )dy (5.35)
−∞ −∞
Let us consider the LHS of the above expression and denote it by
{+∞ [ ] {+∞ [ ]
∗ ∗ d2
ILHS = φ (→y ) Âψ(→y )dy = φ (→y ) − 2 ψ(→y )dy (5.36)
d y
−∞ −∞
RHS of Eq. (5.36) may be integrated using the method of integration of two
multiples
{+∞
d dφ ∗ (y) dψ(y)
ILHS = −φ (→y ) ψ(→y )dy|+∞
∗
−∞ + (5.37)
dy dy dy
−∞
First term in Eq. (5.37) vanishes as both functions φ(y) and ψ(y) are valid state
functions and therefore must vanish at y = ±∞, similarly the derivatives of these
functions must vanish∫at y = ∗±∞.
+∞
Therefore, ILHS = −∞ dφdy(y) dψ(y)
dy
let us integrate this in parts to get
⎧ ⎫ {+∞ 2 ∗ {+∞
dφ ∗ (y) +∞ d φ (y) d2 ∗
ILHS = ψ(y) |−∞ − ψ(y) = − φ (y)ψ(y)
dy dy 2 dy 2
−∞ −∞
∫ +∞ d2 ∗ ∫ +∞ ( d2 φ(y) )∗
or ILHS = −∞ − dy 2 φ (y)ψ(y) = −∞ − dy 2
ψ(y).
Substituting the value of ILHS back in Eq. (5.35), one gets
{+∞ ( 2 )∗ {+∞( )∗
d φ(y)
− ψ(y) = Â †
φ(→
y ) ψ(→y )dy
dy 2
−∞ −∞
( )∗ ( 2 )∗
φ(y)
This give, Â† φ(→y ) = − d dy 2 .
2 2
And Â† = − dy
d
2 ; but − dy 2 = Â.
d
Hence, Â† = Â.

It may be observed that the adjoint Â† of operator Â is equal to the operator.
This is the basic property of a hermitian operator.
SE5.4 Show that parity operator P is Hermitian.
Solution: Parity operator P changes the space coordinates r→ (x, y, z) of the state
−
→
function on which it operates to −r (−x, −y, −z). In order to prove that operator P
is hermitian we have to prove that the adjoint P † is equal to operator P. The adjoint
of the operator is defined by the integral equation
{+∞ [ ] {+∞( )∗
∗
φ (→ r )d x =
r ) Âψ(→ 3
Â† φ(→
r ) ψ(→
r )d3 x
−∞ −∞
We replace operator Â by the parity operator P to get
{+∞ {+∞
∗
[ ] ( † )∗
φ (→ r )d x =
r ) Pψ(→ 3
P φ(→
r ) ψ(→
r )d3 x (5.38)
−∞ −∞
The RHS of the above equation may be denoted by I RHS and may be written as
5.7 Measurement of a Dynamical Variable in Quantum Mechanics 283
{+∞ {+∞
∗
[ ] [ ]
IRHS = φ (→ r )d x =
r ) Pψ(→ 3
φ ∗ (→
r ) ψ(−→
r )d3 x
−∞ −∞
{+∞
[ ]
= −φ ∗ (−→
r ) ψ(→
r )d3 x
−∞
∫ +∞ [ ] ∫ +∞ [ ]
IRHS = −∞ −φ ∗ (−→ r )d3 x = −∞ {Pφ(r )}∗ ψ(→
r ) ψ(→ r )d3 x putting this back
in Eq. (5.38) one gets
{+∞ {+∞
[ ∗
] ( )∗
{Pφ(r )} ψ(→
r )d x = 3
P † φ(→
r ) ψ(→
r )d3 x
−∞ −∞
Comparing the two sides of the above equation gives P = P † which means that
parity operator is hermitian.
Note: If Pψ(r ) = ψ(r ), the wavefunction is said to have even or positive parity, and
if Pψ(r ) = −ψ(r ), the wavefunction is said to have odd or negative parity.
In case Pψ(r ) /= ±ψ(r ), the wavefunction is said to have helicity.
5.7 Measurement of a Dynamical Variable in Quantum

Mechanics
As already mentioned, in quantum mechanics there are operators corresponding to

every observable or measurable quantity, some of which are listed in Table 5.1. Since
measureable quantities are represented by real numbers (and units), operators corre-
sponding to observables are all hermitian operators. The experimentally measured
value of an observable in quantum mechanics is determined by the fourth postulate
of quantum mechanics given below.
Postulate 4 The only measured value of a dynamic variable represented by the
hermitian operator Â, at instant of time ‘t’ when the system is in quantum
mechanical state ψ(→r , t), is one of the eigenvalues of the operator Â.
In quantum mechanics an operator, say, Â may have several eigenvalues, a1 ,
a2 , a3 , … an …. The set of all possible eigenvalues of an operator constitutes the
eigenvalue spectrum. If for each eigenvalue there is a single eigenfunction, the
spectrum of eigenvalues is said to be non-degenerate, i.e.
r , t) = a1 Âϑ1 (→
Âϑ1 (→ r , t); r , t) = a2 Âϑ2 (→
Âϑ2 (→ r , t); r , t) = a3 Âϑ3 (→
Âϑ3 (→ r , t); . . .
r , t) = an Âϑn (→
Âϑn (→ r , t); . . . . . . . . .
(5.39)
The set of all values of ai , i = 1, 2, 3, … n … in expression (5.39) forms a

non-degenerate set or spectrum of eigenvalues.
However, if for some eigenvalue a1 there are more than one eigenfunctions, the
spectrum of eigenvalues is said to be degenerate. If m-different eigenfunctions have
the same eigenvalue a1 , the eigenvalue spectrum is said to be m-fold degenerate.
r , t) = a1 Âϑ1 (→
Âϑ1 (→ r , t); r , t) = a1 Âϑ2 (→
Âϑ2 (→ r , t); r , t) = a1 Âϑ3 (→
Âϑ3 (→ r , t)
(5.40)
Expression (5.40) shows the case of three-fold degeneracy.

The process of measuring a physical variable of the system amounts to operate
the state function of the system by the corresponding operator. Let us consider, as
an example, the case when operator Â (corresponding to some variable ‘a’) has five
non-degenerate eigenvalues ai , i = 1, 2, 3, 4, 5. And the system in consideration is
in state ψ(→r , t) at instant ‘t’. When operator Â will operate on state function ψ(→
r , t),
according to the fourth postulate, one of the five eigenvalues of operator Â will appear
as a result of this operation. Theoretically one may write
r , t) = ai ϑi (→
Âψ(→ r , t) where i may have any value from 1 to 5.
This means that if same measurement is repeated several times, then according to
quantum mechanics, one of the five eigenvalues will appear as the measured value
each time. Quantum mechanics does not tell which eigenvalue value will appear in a
particular measurement. However, quantum mechanics does predict the probability
with which a particular eigenvalue will appear in repeated measurements. Now there
may be three cases: (a) when the eigenvalue spectrum of the operator Â is discrete and
non-degenerate, (b) when the eigenvalue spectrum is discrete but degenerate and (c)
when the eigenvalue spectrum in continuous and not discrete. Method of calculating
probability for a particular eigenvalue in each of the above case is discussed in the
following.
(a) In the case when the eigenvalue spectrum of the operator Â is (discrete
) and
non-degenerate, then quantum mechanics gives the probability P a j that the
measured value will be a particular aj as
|( )|2
( ) | ϑj, ψ |
P aj = (5.41)
(ψ, ψ)
In Eq. (5.41), ϑ j (→
r , t) is the eigenfunction of operator Â corresponding to
eigenvalue aj , i.e.
r , t) = a j ϑ j (→
Âϑ j (→ r , t) (5.41a)
And,
{+∞
( )
ϑj, ψ = r , t)∗ (ψ(→
ϑ j (→ r , t))d3 x (5.42)
−∞
{+∞
(ψ, ψ) = |(ψ(→
r , t))|2 d3 x (5.43)
−∞
In case the state function is normalised to unity, i.e. (ψ, ψ)

∫ +∞ =
−∞
|(ψ(→
r , t))|2 d3 x = 1, probability becomes
| +∞ |2
|{ |
( ) | |
|
P a j = | ϑ j (→ ∗
r , t) (ψ(→ 3 |
r , t))d x | (5.44)
| |
−∞
Important observations: It follows from the above that

(i) Whenever some measurement is done, the process of measurement changes
the state of the system from initial state ψ(→r , t) to a final state ϑ j (→r , t) imme-
diately at the end of the measurement, unless the system was in state ϑ j (→r , t)
when measurement was started (i.e. ψ(→r , t) = ϑ j (→r , t)).
(ii) Quantum mechanics does not tell the exact value of the dynamic variable
that will be obtained in a measurement; it only gives the probability that
a particular value out of the several possible values may be obtained. As
such it will not be wrong to call quantum mechanics as a statistical theory.
(a) In case when the eigenvalue spectrum of operator Â is m-fold degenerate,
the probability P(aj ) is given by
∑m ||∫ +∞ k∗ 3 |
|2
( ) |
k=1 −∞ ϑ j (→
r , t)ψ(→ r , t)d x |
P aj = ∫ +∞ (5.45)
−∞
|(ψ(→ r , t))| d x
2 3
(b) If the operator Â does not have discrete set of eigen values but the eigen-
value spectrum is continuous, then the probability that the measured eigen
value will lie between a, and (a, + da, ) is given by
( ) |ψ(a)|2
dP a , = ∫ +∞ da (5.46)
−∞
|ψ(a , )|2 da ,
SAQ: Under what assumption(s) the wavefunction for a system may be separated
into two independent parts one depending on space coordinates and the other
on time coordinates.
5.7.1 Expectation Value of a Dynamic Variable
Let us once again consider a quantum mechanical system in a state defined by ψ(→ r , t).
Further, let there be a hermitian operator Â that corresponds to a dynamic variable
‘a’. We assume that operator Â has only five discrete non-degenerate eigenvalues
designated as; a1 , a2 , a3 , a4 and a5 with corresponding eigen states φi (→ r , t); i =
1, 2, 3, 4, 5. Now according to the postulate 4 of quantum mechanics, whenever any
attempt will be made to measure the variable ‘a’, the process of measurement will
switch the system from initial state ψ(→ r , t) to one of the eigen states φi (→r , t). The
process of switching will be random; in first attempt the system may be switched
to eigen state φ3 (→
r , t) to give the value of variable ‘a’ as a3 , in the next attempt of
measurement the system may be switched to eigen state φ2 (→ r , t) to give the value as
a2 and so on. Therefore, according to quantum mechanics, the result of measurement
will be one of the eigen values; it may be a1 or a5 , or a3 , or a2 or a4 , only one of the
eigen values and nothing else.
We now consider the actual measurement of the dynamic variable ‘a’ in a labo-
ratory experiment. Suppose the laboratory measurement is repeated three times in
identical conditions and experimental values x 1 , x 2 , x 3 have been recorded. The
normal practice is to take the mean value of these three experimental values as the
final experimental value; so the experimental value of variable ‘a’ is given as
aexp = (x1 + x2 + x3 )/3 (5.47)
Now the question is: To which eigen value of operator Â the experimental value
aexp may be compared? Quantum mechanics says that the experimental value of a
variable should be compared with the expectation value of the variable ‘a’ and not
with any individual eigen value of the operator. The expectation value of a variable
is defined as under.
Expectation value of a variable ‘a’ in state ψ is denoted by
∫ +∞ [ ]
−∞ ψ ∗ (→
r , t) Âψ(→
r , t) d3 x
⟨a⟩ = ∫ +∞ (5.48)
−∞ ψ ∗ (→
r , t)ψ(→
r , t)d3 x
In the case when the state function ψ(→

r , t) is normalised, the expectation value is
given as
{+∞ [ ]
⟨a⟩ = ψ ∗ (→
r , t) Âψ(→
r , t) d3 x (5.49)
−∞
SAQ: Is it possible to measure the value of a dynamical variable when the system
is not in one of the Eigen states of the hermitian operator corresponding to
the variable?
Solved Examples
SE5.5 Calculate the expectation value of variable x 2 for a system in quantum state
ψ(x) = 4e−k(x−b) . Here k and b are real constants.
2
Solution: The expectation value of variable

∫ +∞ ( −k(x−b)2 )∗ 2 ( −k(x−b)2 )
⟨ ⟩ −∞ 4e x 4e dx
x2 = x2 = ∫ +∞ ( 2 )∗ ( 2)
(5.50)
−k(x−b) 4e−k(x−b) dx
−∞ 4e
or
(
∫ +∞ )
2 −2k(x−b)2
⟨ 2⟩ −∞ x e dx
x = ∫ +∞ ( 2)
(5.51)
−2k(x−b) dx
−∞ e
In Eq. (5.51) both the enumerator and the denominator are definite integrals of
standard form with values given as
∫ +∞ ( )
2 −2k(x−b)2 √ π
⟨ 2⟩ −∞ x e dx 3 1
x = ∫ +∞ ( ) = √32k
π
=
−2k(x−b)2 dx 4k
−∞ e 2k
SE5.6 The initial state of a system is defined by the function ψi = 21 ϕ1 + 13 ϕ2 and

ϑn are the normalised eigenfunctions of the total energy operator Ĥ of the
system such that Ĥ ϕn = n 3 K ϕn ; here n is a positive number (n = 1 for state
1 and 2 for state 2) and K a constant having units of energy. What values
for the energy of the system will be obtained and with what probabilities,
in experimental measurements of the system energy? Also calculate the
expectation value of the energy.
Solution: As given in the problem, the initial state of the system is a superposition
state which is the linear combination of two eigen states ϕ1 and ϕ2 of the Hamiltonian
H of the system. Further, according to the fourth postulate of quantum mechanics,
whenever total energy of the system will be measured, the process of measurement
will drive the system to one of the two possible eigen states ϕ1 or ϕ2 .
Suppose, the system is driven to eigen state ϕ1 .
Then using the expression
Ĥ ϕn = n 3 K ϕn ; one gets Ĥ ϕ1 = (1)3 K ϕ1 = K ϕ1 (5.52)

Equation (5.52) tells that the eigenvalue of energy operator in state ϕ1 = K .

Similarly, the eigenvalue of energy operator Ĥ in eigen state ϕ2 = 23 K = 8 K .
Therefore, according to quantum mechanics, measurement of the total energy of
the system will give either K or 8 K as the energy of the system.
Next let us calculate the probability with which the values K and 8 K will occur.
For that we use the relation,
|∫ |2
|( )| | +∞ |
( ) | ϑ j , ψ |2 | −∞ ϑ j (→ r , t)∗ (ψ(→ r , t))| d3 x
P aj = = ∫ +∞
(ψ, ψ) −∞
|(ψ(→ r , t))|2 d3 x
|∫ |2 |∫ ) ||2
| +∞ ∗ | | +∞ ∗ ( 1
| −∞ ϕ1 (ψi )dx | | −∞ ϕ1 2 ϕ1 + 13 ϕ2 dx |
P(a1 = K ) = ∫ +∞ = ∫ +∞ |( )|
−∞
|(ψi )|2 dx | 1 ϕ1 + 1 ϕ2 |2 dx
−∞ 2 3
|∫ |
| +∞ ∗ ( 1 ) ( ∗ 1 )|2
| −∞ ϕ1 2 ϕ1 + ϕ1 3 ϕ2 | dx
= ∫ +∞ |( 1 )|
| ϕ1 + 1 ϕ2 |2 dx
−∞ 2 3
|∫ |
| +∞ ∗ ( 1 ) ( ∗ 1 )|2
| −∞ ϕ1 2 ϕ1 + ϕ1 3 ϕ2 | dx
= ∫ +∞ ( 1 )∗ ( 1 )
−∞ 2 ϕ1 + 3 ϕ2 ϕ + 13 ϕ2 dx
1
2 1
|∫ |
| +∞ ∗ ( 1 ) ( ∗ 1 )|2
| −∞ ϕ1 2 ϕ1 + ϕ1 3 ϕ2 | dx
= ∫ +∞ ( 1 )∗ ( 1 ) ( 1 )∗ ( 1 ) (5.53)
−∞ 2 ϕ1 ϕ + 3 ϕ2 3 ϕ2 dx
2 1
∫ +∞
Equation (5.53) may further reduced using the identity −∞ χm χn = 1, i f m =
n, other wise = 0. |∫ |
| +∞ ∗ ( 1 ) ( ∗ 1 )|2
| −∞ ϕ1 2 ϕ1 + ϕ1 3 ϕ2 | dx
P(a1 = K ) = ∫ +∞ ( 1 )∗ ( 1 ) ( 1 )∗ ( 1 )
−∞ 2 ϕ1 ϕ + 3 ϕ2 3 ϕ2 dx
2 1
Or |∫ | .
| +∞ ∗ ( 1 )|2
ϕ ϕ
| −∞ 1 2 1 | dx
= ∫ +∞ ( 1 )∗ ( 1 ) ( 1 )∗ ( 1 )
−∞ 2 ϕ1 ϕ + 3 ϕ2 3 ϕ2 dx
2 1
Since it is given that eigenfunctions ϕn are normalised,
{+∞ {+∞
| ∗ |2 | ∗ |2
|ϕ ϕ1 | dx = |ϕ ϕ1 | dx = 1
1 1
−∞ −∞
∫ |
1 +∞ | ∗ |2
|
4 −∞
ϕ1 ϕ1 dx 1
P(a1 = K ) = ∫ +∞ | | ∫ | | = 4
= 9/13
1 |ϕ ∗ ϕ1 |2 dx + 1 +∞ |ϕ ∗ ϕ2 |2 dx 1
+ 1
4 −∞ 1 9 −∞ 2 4 9
Similarly, it may be shown that

5.8 Some One-Dimensional Problems 289
∫ +∞ | ∗ |2
| |
−∞ ϕ2 ϕ2 dx
1 1
P(a2 = 8K ) = ∫ +∞ | |
9
∫ | | = 9
= 4/13
1 |ϕ ∗ ϕ1 | dx + 1 +∞ |ϕ ∗ ϕ2 |2 dx
2 1
+ 1
4 −∞ 1 9 −∞ 2 4 9
The expectation value of energy may be calculated using the expression

∫ +∞ [ ] ∫ +∞ ( 1 )∗ [ ( 1 )]
−∞ ψi∗ Ĥ ψi dx −∞ 2 1ϕ + 1
ϕ
3 2
Ĥ ϕ
2 1
+ 1
ϕ
3 2
dx
⟨H ⟩ = ∫ +∞ ∗ = ∫ +∞ ( 1 ) (
∗ 1 )
−∞ ψi [ψi ]dx −∞ 2 ϕ1 + 3 ϕ2 ϕ + 13 ϕ2 dx
1
2 1
∫ +∞ ( 1 )∗ [ ]
ϕ1 + 13 ϕ2 K2 ϕ1 + 8K ϕ dx
3 2
= −∞ 2 13
36
1
K + 89 K 41
= 4
13
= K = 3.15K
36
13
It may be observed that quantum mechanically, the total energy of the system may
assume a value of either 1 K or 8 K. However, the expectation value of energy that
may be compared with experimentally measured energy is 3.15 K.
5.8 Some One-Dimensional Problems
In the following sections we will study the application of quantum mechanics to

some simple problems where the state function of the system will involve only one
spatial direction and the potential energy, if involved in the problem, will be time
independent. Under these conditions one may use one-dimensional Schrodinger’s
time-independent equation (given below) to study the time evolution of the system.
Further, for simplicity, our system will consist of a single particle of mass m.
ℏ2 ∂ 2 x
x )e(− ℏ Et )
i
− φ(→
x ) + V (→ x ) = Eφ(→
x )φ(→ x ); ψ(→
x · t) = φ(→
2m ∂ x 2
5.8.1 Energy States: Bound and Scattering States
Figure 5.1 shows a particle in a one-dimensional potential well. The potential, as

shown in the figure, is specified by the condition that the minimum value of potential
V min occurs at x = x 0 .
The maximum value V 2 at x = x b extends up to positive infinity (+∞); and the
potential takes the value V 1 (< V 2 ) at x = x a . The potential extends up to negative
infinity (−∞) with constant value V 1 . Let us now consider a particle of mass m and
total energy E that is put in the potential well. In classical approach the total energy
Fig. 5.1 Particle in one-dimensional potential
E of the particle may have any value, from E = V min onwards in a continuous way.
However, in quantum mechanics, the particle cannot have all continuous values, the
particle can stay only with some discrete values of energy, say, E 0 , E 1 , E 2 , E 3 , etc.
Particle cannot have energy between E 0 and E 1 ; E 1 and E 2 or E 2 and E 3 and so
on. The energy spectrum of the particle in a potential well is shown in Fig. 5.2. The
lowest energy state with energy E 0 is called the ground state and the state next higher
in energy E 1 the first excited state and so on. The energy spectrum can be divided
into two distinct groups: (i) energy states below V 1 (the lower height of the potential
well) and (ii) states with energy E > V 1 .
Energy states with energies below V 1 are called bound states. It is because if
looked classically, the particle with energy V min < E < V 1 will remain confined
within the region of space defined by the potential well, it will be constrained to
move up to one end of the potential and will then return back at the classical turning
Fig. 5.2 Energy spectrum of a particle in a potential well

point to travel to the other turning point on the other end. Hence these states are
called bound states. Energy states above energy V 1 are called scattering states as
once the particle is in one of the scattering states it is not confined to the potential
well and may scatter away.
Important properties of bound state energy states and corresponding wavefunc-
tions are
1. The bound state energy levels in the case of a one-dimensional potential are
discrete and non-degenerate.
2. The ground state wavefunction φ 0 (x) has no node, which means that in the
space +∞ > x > −∞ the wavefunction does not become zero. Wavefunction
φ1 (x) for the first excited state has one node (becomes zero at one value of x),
wavefunction φ2 (x) for the second excited state has two nodes, φ 3 (x) three
nodes and so on.
SAQ: What will be the difference in the description of bound states looked from
classical physics?
5.8.2 Quantum Mechanical Description of a Free Particle
A particle is said to be a free particle if it has energy E but does not face any potential
V (x).
Space part of one-dimensional time-independent Schrodinger equation describing
a free particle of mass m and energy E may be written as
ℏ2 ∂ 2 x
− φ(→
x ) = Eφ(→
x) (5.54)
2m ∂ x 2
or
∂2x 2m E
φ(→
x ) + 2 φ(→
x) = 0 (5.55)
∂x 2 ℏ
Substituting k 2 = 2m E
ℏ2
; E > 0, above equation reduces to
∂2x
φ(→
x ) + k 2 φ(→
x) = 0 (5.56)
∂x2
Equation (5.56) has two solutions; φ1 = eikx and φ2 = e−ikx that satisfy the
equation. The complete wavefunction ψ(x, t) that contains the time dependence
e(− ℏ Et ) for φ1 and φ2 may be written as
i
ψ1 (x, t) = φ1 e− ℏ t = eikx e− ℏ t = ei(kx− ℏ t )

iE iE E
(5.57)
and
ψ2 (x, t) = φ2 e− ℏ t = e−ikx e− ℏ t = e−i(kx+ ℏ t )

iE iE E
(5.58)
Since ψ1 (x, t) and ψ2 (x, t) are both solutions of Schrodinger equation, it follows
from the principle of quantum mechanical superposition, that their linear combination
ψ(x, t), given below, will also be a solution of Schrodinger equation of the system.
ψ(x, t) = A1 ψ1 (x, t) + A2 ψ2 (x, t) = A1 ei(kx− ℏ t ) + A2 e−i(kx+ ℏ t )

E E
(5.59)
It may be recalled that we made the substitution k 2 = 2m E

2 and that the momentum
ℏ√
p of a particle of mass m and energy E is given as p = 2m E. As such one may
write p = ℏk, with this substitution Eq. (5.59) may be written as
ψ(x, t) = A1 ei(kx− ℏ t ) + A2 e−i(kx+Et) = A1 e ℏ ( px−Et) + A2 e− ℏ ( px+Et)

E i i
(5.60)
First term, A1 e ℏ ( px−Et) , in Eq. (5.60) represents a plane wave travelling in the
i
positive x-direction with well-defined linear momentum p = ℏk and well-defined

energy E = ℏ2mk ; similarly, the second term, A2 e− ℏ ( px+Et) , represents a plane wave
2 2 i
of well-defined linear momentum p = −ℏk and energy E = ℏ2mk moving in negative

2 2
x-direction. It may be observed that there are no boundary conditions and hence no
restrictions on the values of E, which means that the energy of a free particle can
have any value, i.e. energy may have continuous values.
From the derivation of Eq. (5.60) it appears that according to quantum mechanical
treatment, a free particle may be represented as two plane waves moving in positive
x-direction and the other in negative x-direction with well-defined energies (E =
ℏ2 k 2
2m
) and well-defined linear momentum +ℏk and −ℏk, respectively. However, this
is not correct. It is because the first point that must be considered is whether the wave-
function ψ(x, t) given by Eq. (5.60) is a quantum mechanically valid wavefunction
or not. Problems with this wavefunction are
(i) The probability density of finding the particle at some point x for either of the
∫ +∞ φ1 or φ2 is given by
two solutions
P(x) = −∞ φi∗ φi dx is equal to ||Ai |2 (i = 1, 2) which does not depend either
on distance x or on time t. This implies that the total probability of finding
the particle somewhere in the space will tend to become infinite; which is
physically impossible.
(ii) The second problem is that [ the wavefunctions ψ(x, t) could not be normalised
∫ +∞ ∫ +∞ ∫ +∞ ]
because −∞ ψ ∗ ψdx = |A1 |2 −∞ dx + |A2 |2 −∞ dx → ∞, which tends
to infinity.
(iii) Lastly, the velocity with which the particle is moving v is given as v = mp = ℏk
m
while the velocity of the right or left moving plane wave is denoted by vwav =
ℏk
2m
.
It may be observed that according to this description the particle is moving
with twice the velocity of the plane wave which represents it.
It may thus be concluded that plane wave solutions of one-dimensional time-

independent Schrodinger equation cannot be taken as the solution representing a
free particle.
The correct solution of Schrodinger equation for a free particle may be obtained
by superimposing infinite number of coherent plane wave solutions that creates a
wave packet which may be given as
{+∞
1
ψ(x, t) = √ Ak (k)ei(kωx−ωt) dk (5.61)
2
−∞
The amplitude
{+∞
Ak (k) = ψ(x, 0)e−ikωx dx (5.62)
−∞
Here,
E
ω= (5.63)
ℏ
Since an infinite (or sufficiently large) number of plane waves that superimpose on
each other are coherent, they interfere producing a pattern of interference maximums
with largest amplitude at x = 0 as shown in Fig. 5.3. The intensity of successive
maximums decreases with the increase of distance x on both sides of the origin.
The envelope of the interference maximum makes the wave packet that is localised
around x = 0.
Fig. 5.3 Snapshot of the wave packet formed by the interference of large number of plane waves.
The wave packet is localised at x = 0 and moves in +x-direction with the group velocity V group
Since the exponential factor eikx oscillates rapidly, the wavefunction ψ(x, t) given
by Eq. (5.61) undergoes destructive interference and vanishes at x = ∓∞. In a
similar way, the amplitude function Ak (k) defined by Eq. (5.62) is localised at k =
0 in k-space and goes to zero for large values of k.
The size of the wave packet in space is specified by the half-width ∓Δx where the
intensity of maximum [|ψ(x = ∓Δx)|2 ] drops to √1e of its maximum value at origin
[|ψ(x = 0)|2 ]. Similarly, the size of the wave packet corresponding to the amplitude
function Ak (k) in k-space is defined by the half-width ∓Δk, such that the magnitude
of the function [|Ak (k = ∓Δx)|2 ] at k = ∓Δk is √1e of [|Ak (x = ∓Δx)|2 ]. That is
[|ψ(x=∓Δx)|2 ] = [|Ak (k=∓Δx)|2 ] = √1 .
[|ψ(x=0)|2 ] [|Ak (k=0)|2 ] e
It can be shown that both ψ(x) and Ak (k) are normalised to unity.
One may ask the physical interpretation of the wave packet. The physical inter-
pretation of the wave packet may be given as follows; |ψ(x, t)|2 is the probability
density that the particle is found at point x at time t, while |ψ(x, t)|2 dx may be
interpreted as the probability of finding the particle in the interval x and (x + dx).
Similarly, |Ak (k)|2 gives the probability density of measuring the wave vector value
as k or momentum value p = k/ℏ of the particle. Also, |Ak (k)|2 dk may be defined as
the probability of measuring the value of wave vector between k and (k + dk).
Identifying a free particle by a wave packet removes all difficulties associated with
plane wave representation of the free particle. Wave packet and the wavefunctions
are normalised; wave packet description provides probability density which never
approaches infinity, momentum and location of the particle are not known exactly at
a given time. Problem regarding the speed of the particle and the plane wave is no
more, particle which is represented by the wave packet moves with the group velocity
while individual waves travel with the phase velocity. Hence there is no confusion
about the velocity of the particle and the velocity of the wave.
SAQ: In quantum physics a free particle is represented by wave packet. What is the
mechanism by which the wave packet is formed?
5.8.3 Particle in a One-Dimensional Asymmetric Infinite

Potential Well
A one-dimensional infinite potential well of small width ‘a’ is shown in Fig. 5.4.
The potential well is asymmetric because it is not symmetrical about the vertical axis
through x = 0. The potential V (x) is zero between x = 0 and x = a; and for all other
values of x it is infinite. In mathematical language the potential well may be defined
as
⎧
0, for 0 < x < a
V (x) = (5.64)
∞, for x ≤ 0, x ≥ a
This is also called the problem of a particle in a box.

Fig. 5.4 One-dimensional infinite potential well of small width ‘a’
We wish to study the energy states for a particle of mass m and energy E placed
inside the potential well. Let us first look the problem from the point of view of
classical physics. Since the particle is confined from both sides by infinite potential
walls, the particle will not be able to cross the potential walls at x = 0 and at x = ∞,
will turn back
√
when it hits the potential walls and will move with constant speed v
= p/m = ∓ 2m m
E
from one end to the other end of the well. From the viewpoint of
quantum physics, since the particle is confined in a small space it corresponds to the
problem of bound states.
The time-independent Schrodinger equation for this case (as in the case of a free
particle) may be written as
ℏ2 ∂ 2 x
− φ(→
x ) = Eφ(→
x)
2m ∂ x 2
or
∂2x 2m E
φ(→
x ) + k 2 φ(→
x ) = 0, where k 2 = 2 ; E > 0 (5.65)
∂x 2 ℏ
Both φ1 (→ x ) = e−ikx satisfy the second-order partial differential

x ) = eikx and φ2 (→
equation (5.65).
The linear combination of these functions with arbitrary (may be complex)
constants A and B may be written as
ϑ(x) = Aφ1 (→ x ) = Aeikx + Be−ikx

x ) + Bφ2 (→ (5.66)
Function ϑ(x) will also be a solution of Eq. (5.65). It is now required to determine
the value of arbitrary constants A and B. To find the values of these constants, we
review the boundary conditions; at x = 0 and x = a, where the potential is infinite,
function ϑ(x) must vanish.
Therefore, when
x = 0, ϑ(0) = A + B = 0 and B = −A (5.67)
Also, when
x = a, ϑ(a) = Aeika + Be−ika = 0
One may expend exponential terms in sine and cosine forms to get
or
A{cos(ka) + i sin(ka)} + (−A){cos(ka) − i sin(ka)} = 0
or
2iA sin ka = 0 (5.68)
In Eq. (5.68), arbitrary constant A is not zero; therefore, sin ka = 0; this is possible
only when
ka = nπ, where n may have integer values n = 1, 2, 3, . . . (5.69)
The value n = 0 is not taken because in that case the function ϑ(x) will become
zero everywhere. √
On substituting the value of k = 2mℏ
E
in Eq. (5.69) one gets
√
2m E π 2 ℏ2
a = nπ ; or E = n 2 ; n = 1, 2, 3 . . . (5.70)
ℏ 2ma 2
Equation (5.70) tells that a particle confined in an infinite potential well may have
many discrete values for its energy E corresponding to different values ( 2 of2 )integer n.
π ℏ
The lowest energy state, called the ground state, has energy E 0 = 12 2ma 2 ; the first
( 2 2)
π ℏ
excited state has the energy E 1 = (2)2 2ma 2 and so on. It may be observed that the
energy separation between successive energy states is not uniform; it increases with
the increase of the value of n. Following conclusions may be drawn from Eq. (5.70):
(i) Particle in a box can move within the infinite potential well only with fixed
π 2 ℏ2
discrete value of energies given by, E (n−1) = n 2 2ma 2 , where n can have integer
values, 1, 2, 3, …. This is in contrast to the classical picture in which particle

in a box may oscillate within the well with any value of energy starting from
zero. In quantum mechanical picture the minimum energy that a particle can
π 2 ℏ2
have is 2ma 2.
Fig. 5.5 Spectrum of energy states for a particle in a one-dimensional box
(ii) There is infinite number of discrete energy levels or states that a particle in an
infinite potential well may have. Separation between two consecutive energy
states increases with the value of the positive nonzero integer n.
(iii) Energy E (n−1) is inversely proportional to the square of the width of the potential
well; i.e. E (n−1) ∝ a12 .
It may be noted that energy E (n−1) is the kinetic energy of the particle, since the
potential energy in the region between two infinite potential walls is zero. Figure 5.5
shows the spectrum of allowed energy states of a particle in an infinite potential well.
As shown in this figure, the minimum energy that a particle may have inside the box
is E 0 and the higher energy states are not equidistant, rather the separation increases
with the energy of the state. The lowest energy state E 0 is called the ground state and
the next higher energy state E 1 as the first excited state, E 2 the second excited state
and so on.
5.8.3.1 Eigenfunction for a Particle in a One-Dimensional Infinite Box
Since the energy eigenvalue E (n−1) depends on the value of the positive integer n,
there will be n eigenfunctions corresponding to each energy value. The space part of
eigenfunctions, using Eq. (5.68), may be written as
( nπ )
ϑ(n−1) (x) = Dn sin x (5.71)
a
where Dn is a constant, the value of which may be obtained using the condition of
normalisation given below
{+∞ {+a
∗ ∗
|Dn | 2
ϑ(n−1) ϑ(n−1) dx = |Dn | 2
ϑ(n−1) ϑ(n−1) dx = 1
−∞ 0
or
{+a {+a( ( nπ ))2

∗
|Dn | 2
ϑ(n−1) ϑ(n−1) dx = |Dn | 2
sin x dx = 1
a
0 0
∫a
We now make the use of the standard result 0 sin2 (θ ) = a
2
− sin42θ |a which gives
{a } /
2
|Dn |2
= 1 or Dn = (5.72)
2 a
Therefore, the normalised space part of the eigenfunction for (n − 1) energy state
may be written as
/
2 ( nπ )
ϑ(n−1) (x) = sin x (5.73)
a a
Eigenfunctions, ϑ0 (x), ϑ1 (x) and ϑ2 (x) for the ground state, first excited state
and the second excited states within the width ‘a’ of the potential well are shown
in Fig. 5.6. It may be observed that ϑ0 (x) does not have zero value for any value of
x, it does not have any node within the width of the potential well, ϑ1 (x) has one
node as it gets a zero value for x = 0.5a. ϑ3 (x), and the normalised space part of the
eigenfunction for the second excited state has two nodes, indicated by arrows in the
figure. In general it may be said that the eigenfunction for the (n − 1)th excited state
will contain n-nodes.
The total wavefunction including the time-dependant part may be written as
ψ(n−1) (x, t) = space part × time dependent part = ϑ(n−1) (x)e−i( ℏ )t

En
or
/
2 ( nπ ) −i
( ) ( )
n2 π 2 ℏ n2 π 2 ℏ
−i t t
ψ(n−1) (x, t) = ϑn (x)e 2ma 2 = sin x e 2ma 2 (5.74)
a a
Since Schrodinger equation is linear, the most general stationary state solu-
tion for one-dimensional infinite well is given as the linear combination of several
wavefunctions of type (5.74) with different multiplicative constants Cn
Fig. 5.6 Space part of eigenfunctions for ground, first and second excited states for a particle in
infinite potential
/
2 ( nπ ) −i
∞ ( )
gen
∑ n2 π 2 ℏ
t
ψ(n−1) (x, t) = Cn sin x e 2ma 2 (5.75)
n=1
a a
5.8.3.2 Extension to Two-Dimensional and Three-Dimensional Infinite

Potentials
Results obtained for one-dimensional case may be extended to the two and the three-
dimensional infinite potential wells of width ‘a’ in each direction. In the case of 1D
box the energy levels of the particle are specified by Eq. (5.70) as given below
π 2 ℏ2
E (n−1) = n 2
2ma 2
In a two-dimensional case, the corresponding expression for energy may be written
as ( )( π 2 ℏ2 )
E (n−1) = n 2x + n 2y 2ma 2 where nx and ny may have values 1, 2, 3, …
So the ground state energy of the particle in a two-dimensional well of width ‘a’
will be
( ) ( 2 2)
( ) π 2 ℏ2 π ℏ
E 0,0 = 12 + 12 2
= 2
2ma 2ma 2
( )( π 2 ℏ2 ) ( 2 2)
π ℏ
And the first excited state energy as E 0,1 = 12 + 22 2ma 2 = 5 2ma 2 and so
on.
In the three-dimensional case the energy may be written as(
( )( π 2 ℏ2 ) π 2 ℏ2
)
Ground state energy E (0,0,0) = 12 + 12 + 12 2ma 2 = 3 2 .
( 2 ) ( 2 2 )2ma ( 2 2 )
π ℏ π ℏ
Energy of first excited state E (0,0,1) = 1 + 12 + 22 2ma 2 = 6 2ma 2 .
5.8.4 Potential Barrier and Tunnelling
In the last section we studied the problem of a particle in an infinite potential well.
In this section we will study the quantum mechanical behaviour of a particle of mass
m and kinetic energy E which is projected on a one-dimensional potential barrier of
height V 0 and of width ‘a’. A potential barrier is a space of width ‘a’ where there is
a potential of constant magnitude V 0 all over the space, and this space with potential
V 0 is surrounded from all sides by the free space where there is no potential, as
shown in Fig. 5.6. The potential barrier may be defined as
⎧
⎨ = 0 when X < 0
V = = V0 when 0 ≤ X ≤ a
⎩
= 0 when X > a
The incident particle, which in quantum mechanical language is a wave packet,

meets different fates depending on the energy E of the particle, the height V 0 and
the width ‘a’ of the potential barrier. Three distinct cases may happen. First, when
the height of the potential barrier V 0 is infinite and the width of the barrier ‘a’ is
also infinite; in this case the incident wave packet of the particle bounces back, not
being able to penetrate the barrier. This is just like the case of classical physics where
the incident particle is reflected back by the infinite potential. In the second case the
barrier height V 0 is finite and the width is infinite, in this case it can be shown, that the
incident particle or the wave packet penetrates into the barrier, travels some distance
in the region of the barrier and then dies out. However, in case when the barrier
height V 0 is finite and larger than the energy E of the particle and the width of the
barrier ‘a’ is small, the particle may penetrate the barrier and travel through the small
width of the barrier to appear on the other side of the barrier. This penetration of the
small width finite barrier by the incident wave packet of lower energy (representing
the incident particle) is called barrier tunnelling. There is no analogue of barrier
tunnelling in classical physics. We will show in the following how using quantum
mechanics the phenomena of barrier tunnelling may be explained (Fig. 5.7).
In the framework of quantum mechanics, the present problem where the one-
dimensional potential barrier does not depend on time, may be treated using time-
independent one-dimensional Schrödinger equation given as
Fig. 5.7 Particle of mass m and energy E impinging on a one-dimensional potential barrier of
height V 0 and width a
ℏ2 d2 φ(x)
− + V (x)φ(x) = Eφ(x) where − ∞ < x < +∞ (5.76)
2m dx 2
Here, φ(x) represents the wavefunction of the particle.
Since the interesting case of barrier tunnelling occurs when the energy of the
particle is less than the barrier height, we assume that V 0 > E. Further, the total
one-dimensional space may be divided in to three regions: Region-I from x = −∞
to x = 0 where V (x) = 0; Region-II from x = 0 to x = a where V (x) = V 0
and Region-III from x = a to x = +∞ where V (x) = 0. If ϑ1 (x), ϑ2 (x) and ϑ3 (x),
respectively, denote the space part of the particle wavefunctions in Region-I, Region-
II and Region-III, then Schrodinger equations for the three regions may be written
as
Region-I
ℏ2 d2 ϑ1 (x) d2 ϑ1 (x) 2m E
− = Eϑ1 (x) ⇒ + kI2 ϑ1 (x) = 0, where kI2 =
2m dx 2 dx 2 ℏ2
(5.77)
Region-II
ℏ2 d2 ϑ2 (x) d2 ϑ2 (x)
− + V0 ϑ2 (x) = Eϑ2 (x) ⇒ = kII2 ϑ2 (x);
2m dx 2 dx 2
2m(V0 − E)
where kII2 = (5.78)
ℏ2
Region-III
ℏ2 d2 ϑ3 (x) d2 ϑ3 (x) 2m E
− = Eϑ3 (x) ⇒ + kI2 ϑ3 (x) = 0, where kI2 =
2m dx 2 dx 2 ℏ2
(5.79)
It may be noted that Region-I and Region-III have no potential and we are consid-
ering a particle
/ of mass M and energy E; therefore, the value of the wave number
kI = kIII = 2Mℏ2
E
.
We now attempt to write wavefunctions for the three regions. Let us first consider
Regions-I and III where the potential is zero and regions are free.
Region-I
Let us examine Eq. (5.77). Wavefunction ϑ1 (x) = Ae+ikI x + Be−ikI x , where A and
B are arbitrary complex constants, satisfies Eq. (5.77), which may be verified by
differentiating the wavefunction ϑ1 (x) twice with respect to x. As may be observed,
Ae+ikI x represents a plane wave moving in positive x-direction and Be−ikI x represents
a plane wave moving in negative x-direction. It may be noted that wave represented
by Aeikx = A(cos kx + i sin kx) or by Be−ikx = B(cos kx − i sin kx) has oscillatory
nature. What happens is that the incident wave when hits the potential barrier at x
= a, a part of the wave gets reflected back. The reflected wave moving in negative
x-direction is represented by ϑIref = Be−ikI x . Hence, one may write
ϑ1 (x) = Ae+ikI x + Be−ikI x (5.80)
where reflected wave is given by
ϑIref = Be−ikI x (5.81)
Region-III
Similarly, ϑ3 (x) = Fe+ikI x + Ge−kI x satisfies Eq. (5.79) for Region-III. However, it
may be realised that there must not be any reflected wave in Region-III as the potential
free space extends up to +∞. This means that constant G in above expression must
be zero. Hence the wavefunction in Region-III may be written as
ϑ3 (x) = Fe+ikI x (5.82)
Region-II
Wavefunction ϑ2 (x) must satisfy Schrodinger Eq. (5.78)
d2 ϑ2 (x) 2m(V0 − E)
= kII2 ϑ2 (x); where kII2 =
dx 2 ℏ2
Since it is assumed that V 0 > E, kII2 is a real positive quantity and kII will also be
real quantity. The general solution to the Schrodinger equation in Region-II is given
by
ϑ2 = CekII x + De−kII x (5.83)
It is worth noting that the wavefunction in Region-II shows two waves, one De−kII x
that is not oscillatory but decays exponentially with x and the other CekII x that
increases exponentially with x. The wavefunction in Region-III may thus be written
as
ϑ2 = CekII x + De−kII x (5.84)
5.8.4.1 Boundary Conditions
Having derived expressions for the space part of wavefunctions in three regions,
we now proceed to find the value of unknown constants, A, B, C, D and F. To
obtain the values of these constants, one uses the boundary conditions that the three
wavefunctions must obey.
(i) Condition of continuity demands that wavefunctions on the two sides of the
potential boundary should match each other, that means
ϑ1 (x) = ϑ2 (x) at x = 0 and ϑ2 (x) at (x = a) = ϑ3 (x) (at x = a)
Putting
ϑ1 (0) = ϑ2 (0) gives A + B = C + D (5.85)
Putting
ϑ2 (a) = ϑ3 (a) gives CekII a + De−kII a = FeikI a (5.86)
(ii) Condition of smooth joining demands that the first derivatives of wavefunc-
tions on the two sides of the potential boundary evaluated at the boundary must
be equal.
dϑ1 (x) dϑ2 (x) ikII

= |x=0 or (A − B) = − (C − D) (5.87)
dx dx kI
And
dϑ2 (x) dϑ3 (x) ikI ikI a
= |x=a that gives CekII a − De−kII a = Fe (5.88)
dx dx kII
It follows from Eqs. (5.86) and (5.88), once by adding and by subtracting, that
( )
F ikI a ikI −kII a
C= e 1+ e (5.89)
2 kII
and
( )
F ikI a ikI kII a
D= e 1− e (5.90)
2 kII
Substitution of these values of C and D in Eqs. (5.85) and (5.87) gives

⎧( ) ( ) ⎫
B F ikI a ikI −kII a ikI kII a
1+ = e 1+ e + 1− e
A 2A kII kII
⎧ ( ka )⎫
F ikI a e + e
kII a −kII a
ikII e − e−kII a
II
= e −
A 2 kI 2
⎧ ⎫
F ikI a ikII
= e cosh(kII a) − sinh(kII a) (5.91)
A kI
and
⎧ ⎫
B F ikI a ikI
1− = e cosh(kII a) + sinh(kII a) (5.92)
A A kII
Adding Eqs. (5.91) and (5.92) one gets

⎧ ( ) ⎫
F ikI a kII kI
2= e 2 cosh(kII a) + i − sinh(kII a) (5.93)
A kI kII
and
F ikI a 2
e ={ ( ) } (5.93a)
A 2 cosh(kII a) + i kI −
kII kI
sinh(kII a)
kII
Also by subtracting Eq. (5.92) from Eq. (5.91) one gets

( )
B F ikI a kII kI
2 = −i e − {sinh(kII a)} (5.94)
A A kI kII
F ikI a
We substitute the value of A
e from Eq. (5.93a) in Eq. (5.94) to get
( )
B
kII
kI
{sinh(kII a)}
− kI
kII
= −i { ( ) } (5.95)
A 2 cosh(kII a) + i kkIII − kkIII sinh(kII a)
Also, from Eq. (5.93a)
F 2e−ikI a
={ ( ) } (5.96)
A 2 cosh(kII a) + i kkIII − kI
sinh(kII a)
kII
|F|2 |B|2
From Eqs. (5.95) and (5.96) one may calculate |A|2
and |A|2
as given below
( )∗ ( )
|F|2 F∗ F F F
= ∗ =
|A| 2 A A A A
⎛ ⎞∗
−ikI a
2e
= ⎝{ ( ) }⎠
2 cosh(kII a) + i kI − kII sinh(kII a)
kII kI
⎛ ⎞
−ikI a
⎝{ 2e
( ) }⎠
2 cosh(kII a) + i kkIII − kkIII sinh(kII a)
⎛ ⎞
|F|2 ⎜ 4 ⎟
=⎝ ( 2 2 )2 ⎠ (5.97)
|A|2 k −k
4 cosh2 (kII a) + kIIII K II sinh2 (kII a)
and
( 2 2 )2
kII −kI
|B|2 kII K I
sinh2 (kII a)
= ( ( 2 2 )2 ) (5.98)
|A|2 kII −kI
4 cosh (kII a) + kII K I
2
sinh (kII a)
2
5.8.4.2 Transmission Coefficient
If v denotes the velocity of the incident particle, then the flux incident on the potential
barrier is given by
Flux(inci) = vϑ1∗ (x)ϑ1 (x) = v A∗ A = v|A|2 (5.99)

Flux transmitted across the barrier
Flux(trans) = vϑ3∗ (x)ϑ3 (x) = v|F|2 (5.100)
Transmission coefficient T , which is the probability that the incident particle will
cross through the potential barrier and appear on the other side of it, is defined as the
ratio of the transmitted flux to the incident flu and may be given by
Flux(trans) v|F|2 |F|2

T = = =
Flux(inci) v|A|2 |A|2
⎛ ⎞
⎜ 4 ⎟
=⎝ ( )2 ⎠ (5.101)
kII2 −kI2
4 cosh (kII a) +
2
kII K I
sinh (kII a)
2
Equation (5.101) tells that transmission coefficient has a finite nonzero value indi-
cating that a quantum particle has a finite probability of going across a potential barrier
of height larger than particle’s energy, a phenomena which is not possible in classical
physics. The process of going across the potential barrier of height V 0 by a quantum
particle of energy E, where V 0 > E, is called quantum mechanical tunnelling. Many
phenomena in physics, like alpha radioactive decay, transfer of charge in digital elec-
tronics, etc., can be explained only on the basis of quantum tunnelling. Tunnelling
is a quantum mechanical effect without and classical counterpart.
Equation (5.101) for the transmission coefficient may be rewritten using the
identity cosh2 (kII a) = 1 + sinh2 (kII a) as
1
T =⎧ ( 2 2 )2 ⎫ (5.102)
k +k
1 + 41 kIIII K II sinh2 (kII a)
We now substitute the values of kI and kII in above equation from Eqs. (5.78) and
(5.79) as
kI2 = 2m
ℏ2
E
and kII2 = 2m(Vℏo2−E) to get
1
T ={ ( √ )} (5.103)
V02
1+ 1
4 E(V0 −E)
sinh 2 a
ℏ
2m(V0 − E)
In the case when the barrier height V 0 is much larger than energy E of the particle,
i.e.
/
(a √ ) a √2mV E
0
2m(V0 − E) = 1− ≫1 (5.104)
ℏ ℏ V0
And therefore, one may write

(a √ ) 1 a√2m(V0 −E )
sinh 2m(V0 − E) ≈ e ℏ (5.105)
ℏ 2
( √ )
Substituting this value for sinh ℏa 2m(V0 − E) in Eq. (5.103) one gets
1 1
T ={ ( √ )} = ⎧ √ ⎫
V02 2m (V0 −E )
1+ 1
4 E(V0 −E)
sinh2 ℏa 2m(V0 − E) 1+ 1 V02 1
e
2a
ℏ
4 E(V0 −E) 4
or
( ) ( √ )
16E E − 2a 2mℏ(V0 −E )
T = 1− e (5.106)
V0 V0
It may be observed that Eq. (5.106) gives the magnitude of the transmission
coefficient in the low energy limit. It may further be shown that in case when E and
V 0 are comparable, i.e. E ∼ V0 , transmission coefficient becomes
1
T =( ) (5.107)
ma 2 V0
1+ 2ℏ2
5.8.4.3 Reflection Coefficient
As mentioned earlier, when the incident particle beam hits the potential barrier at x
= 0, a part of the incident beam gets reflected. The reflection coefficient denoted as
R and defined as the ratio of the reflected flux to the incident flux may be written as
⎛ ( )2 ⎞
kII2 +kI2
|B| 2 sinh 2
(k a)
⎜ II
kII K I ⎟
R= =⎝ ( 2 2 )2 ⎠ (5.108)
|A|2
kII −kI
4 cosh (kII a) + kII K I
2
sinh (kII a)
2
In the limit E ∼ v0 reflection coefficient reduces to
1
R= ( ) (5.109)
2ℏ2
1+ ma 2 V0
Tunnelling and reflection are quantum phenomena based on the wave nature of
matter, when a wave hits the boundary separating the two media a part of the wave is
reflected and a part is transmitted. In case the width of the second medium is small,
the transmitted wave propagates to the other end of the second medium and may be
transmitted across the second boundary.
Fig. 5.8 Matching and

smooth joining of
wavefunctions at the
boundaries of the potential
barrier
It may be observed that the part of the wave transmitted in the potential barrier does
not have oscillatory nature, rather it decays exponentially. Matched and smoothly
joined wavefunctions at the boundaries of the potential barrier are shown in Fig. 5.8.
SAQ: It has been shown that a free particle may be represented by a wave packet
and not by a plane wave. However in this example of quantum tunnelling the
wavefunction for the free particle is taken as a plane wave. Can you give a
reason why is it justified?
5.9 Heisenberg Uncertainty Principle
We have seen how quantum mechanics predicts non-uniform discrete energy levels
for a particle confined in a limited space, in contrast to the continuous energy states
predicted by classical physics. Similarly, quantum mechanical tunnelling has no
parallel in classical physics. Uncertainty principle, put forward by Heisenberg in
1927, is another characteristic of quantum mechanics. It tells that there is fuzziness
in nature, particularly at microscopic level. In order to appreciate the uncertainty
principle it is required to define pairs of dynamic observables called canonical
conjugates. Typical canonical conjugates are position of a particle r→ and its linear
momentum − →p ; their three components; position coordinate x and x-component of
linear momentum px ; position coordinate y and y-component of linear momentum
p y and position coordinate z and z-component of linear momentum pz form pairs
of conjugate variables. Another pair of canonical conjugates is the energy E of the
particle and time t at which the energy is measured. According to this principle, each
observable of the canonical pair cannot be measured simultaneously with complete
accuracy. According to Heisenberg uncertainty principle, it is not possible to nail
down with absolute accuracy the speed of a microscopic particle and its location
simultaneously, more one tries to accurately determine the speed less he knows
about the position of the particle. If Δx and Δpx are, respectively, the uncertainties
in the measurement of the x-coordinate of a particle and its linear momentum px in
x-direction at some instant of time, then according to the uncertainty principle
Δx · Δpx ≥ ℏ/2
5.10 Correspondence Principle and Ehrenfest’s Theorem 309
Similarly,
Δy · Δp y ≥ ℏ/2; Δz · Δpz ≥ ℏ/2; and ΔE · Δt ≥ ℏ/2 (5.110)
Here, ΔE and Δt are the uncertainties in the measurement of time t and energy E
of the particle. A common misconception about the principle of uncertainty is that it
arises because of the limitations of measuring instruments; however, the fact is that
in spite of the availability of most accurate and precise measuring instruments, the
law puts limits on the minimum uncertainties in measured values. The argument put
forward in support of the uncertainty principle is that the very process of measurement
alters the conditions of the system introducing uncertainties. Though Heisenberg
derived the principle, however, the derivation is beyond the scope of our present
discussion.
Uncertainty principle has several applications in quantum physics; for example the
excited states of a system have a certain mean life, say, Δτ , which is the uncertainty
in time, and therefore, the uncertainty ΔE in the energy of the excited state, from
uncertainty principle may be given by
ℏ
ΔE ≈ (5.111)
2Δt
where ΔE is called the width of the excited state. It may be noted that in most
cases the ground state of the system is very stable which means that Δt is large, and
therefore, the energy width of the ground state is small, i.e. ground states are sharp.
SAQ: What may be the probable reason for uncertainty in quantum mechanical
measurements?
5.10 Correspondence Principle and Ehrenfest’s Theorem
Correspondence principle essentially demands that any new theory, under suitable
approximations/extensions must merge or dissolve into older theories. Particularly,
in the case of quantum mechanics which works for microscopic systems and is
characterised by quantum numbers, in the limit of large quantum numbers it should
give same results as are predicted by classical mechanics.
Bohr was the first scientist who gave his well-known thesis on Complementarity
and its Copenhagen Interpretation. In fact Bohr put three different aspects of his corre-
spondence principle: first is the frequency interpretation, according to which corre-
spondence principle is a statistical asymptotic agreement between one component
in the Fourier decomposition of the classical frequency and the quantum frequency
in the limit of large quantum numbers. Secondly, there is intensity interpretation
according to which it is a statistical agreement in the limit of large quantum numbers
between the quantum intensity and the classical intensity. Quantum intensity may
be understood in terms of the probability of a quantum transition while classical
intensity as the square of the amplitude of one component of classical motion. The
third interpretation deals with selection rules, according to which the correspondence
principle is the statement that each allowed quantum transition between stationary
states corresponds to one harmonic component of the classical motion.
Ehrenfest’s theorem, in a way, also establishes a bridge between the quantum
mechanics and the classical physics. The theorem states that ‘The average values of
observables in quantum mechanics obey the classical mechanics’. As an example,
it is possible to write the equation of motion for the expectation values of the position
momentum operators, given below, exactly in the same way as the equation of motion
in classical physics.
( ( )) / \
d2 x̂ ∂ V (x)
= − (5.112)
dx 2 ∂x
Equation (5.112) could be derived using tools of quantum mechanics, but it looks
like Newton’s equation of motion where force and hence acceleration may be derived
from the gradient of the potential.
SAQ: What is the logic behind Correspondence principle?
Solved Examples
SE5.7 Calculate the energy difference between the 2nd excited state and the ground
state of a 1D, 2D and 3D infinite potential wells of same width a in each
direction.
π ℏ2 2
Solution: Let us denote 2ma 2 = ε0 .
The energy difference between the second excited state and the ground state
(E 2 − E 0 )(1-D) = (9ε0 − 1ε0 ) = 8ε0
(E 2 − E 0 )(2-D) = (8ε0 − 2ε0 ) = 6ε0
(E 2 − E 0 )(3-D) = (9ε0 − 3ε0 ) = 6ε0
SE5.8 A particle of mass m is confined in a 1D potential well of width b, if the

normalised wavefunction ψ(x) for the particle in the range x = 0 to x = b
is given as
( )
ψ(x) = Bx b2 − 2x 2
and is zero for all other values of x, calculate the value of constant B.
Solution: It is given that the wavefunction is normalised, hence

{+∞
ψ ∗ ψdx = 1
−∞
∫ +∞ ∗ ∫b ∗ ∫ b[ ( 2 )]
2 2
But −∞ ψ ψdx = 0 ψ ψdx = 0 Bx b − 2x dx =
∫ [ {
2 b 2 4
}]
B 0 x b − 4b x + 4x dx.
2 2 4
∫b[ { }] ∫ b [{ }]
Or B 2 [ 0 x 2 b4 − 4b2 x 2 +] 4x 4 dx = B 2 0 b4 x 2 − 4b2 x 4 + 4x 6 dx = 1.
3 2 5 7
Or B 2 b4 x3 − 4b5x + 4x7 |b|0 = 1.
{ 7 7 7
} ( 7) /
Or B 2 b3 − 4b5 + 4b7 = 1 Or B 2 11b 105
= 1. So, B = 105
11b7
.
SE5.9 In a particular molecule the position of a certain ion of mass 6.0 × 10–26 kg
can be determined to an accuracy of 1 µm. Calculate the accuracy with which
the speed of the ion may be determined.
Solution: It is given that the uncertainty in the position of the ion Δx = 1 µm =

1 × 10−6 m; mass M of the ion is given as 6.0 × 10–26 kg; value of ℏ = (6.626 ×
10−34 /2 × 3.14) J s.
If the uncertainty in the velocity of the ion is Δv; then Δpx = MΔv = 6.0 ×
10−26 × Δv.
However, from uncertainty principle: Δx ·Δpx = 1×10−6 ×6.0 ×10−26 ×Δv ≈
(6.626 × 10−34 /2 × 3.14).
6.626×10−34 −4
Or Δv ≥ 2×2×3.14×1×10 −6 ×6.0×10−26 ≈ 8.8 × 10 m/s.
SE5.10 Calculate the probability that an electron of mass 9.1 × 10−31 kg and of
energy 1.0 eV may tunnel through a potential barrier of height 1.5 eV and
width 0.5 nm.
Solution: First of all let us convert all given quantities in SI units.

Energy of electron E = 1 eV = 1 × 1.6 × 10−19 J; barrier height V0 = 1.5 ×
1.6 × 10−19 J.
Barrier width a = 0.5 × 10−9 m. And ℏ = 1.055 × 10−34 J s.
We use the ( √ expression) of the transmission coefficient T =
( ) − 2a 2m(V0 −E )
ℏ
16E
V0
1 − VE0 e
( √ )
(
−19 ) − 2×0.5×10−9 (
2×9.1×10−31 0.5×1.6×10−19 )
16 × 1.6 × 10 1 1.055×10−34
T = 1− e
1.5 × 1.6 × 10−19 1.5
[√ ] [ −9 ×3.816×10−25
]
8 −1×10−9 14.56×10−50 − 1×101.055×10
or T = 2.25
e = 3.56e
1.055×10−34
= 3.56 × e−3.61 =
−34
3.56 × 0.027.
Hence transmission probability T = 0.096 = 9.6 × 10−2 .
Problems
P5.1 Calculate the transmission coefficient for an electron of mass 9.1 × 10–31 kg
and kinetic energy 1.3 eV for a potential barrier of height 1.5 eV and width
0.5 nm. The value of rationalised Planck constant ℏ is 1.055 × 10–34 J s.
ANS: 0.187
d2 x
P5.2 Calculate the eigenvalue of function sin(2x) for operator Â = dx2
.
ANS: −4
P5.3 A particle of mass m is confined in a one-dimensional infinite potential well
of width a. The particle at instant t is in state χ (x) characterised by χ (x) =
B{θ1 (x) + θ2 (x)}, where θ1 (x) and θ2 (x) are, respectively, the normalised
wavefunctions for the ground and third excited states of the system and B
the normalisation constant. Determine the value of B and the average value
Δ
of operators x̂ and px .
⟨ ⟩ ⟨ ⟩
Δ
ANS: B = √12 , x̂ = 2a , px = 0
P5.4 The mean life of an excited nuclear state is 4 ns, what will be the order of its
energy width?
ANS: 8.2 × 10–8 eV
P5.5 A metallic ball of 2.0 g is constrained to roll inside a glass tube of 2.0 cm
length, which is closed at both the ends. If this ball is considered similar to
a particle moving in one-dimensional infinite square well, what is the value
of quantum number n if the ball is initially given an energy of 2.0 × 10−3 J?
ANS: n = 1.77 × 1029
P5.6 Calcuate the de Broglie wavelength of a proton moving with a speed of
3 × 105 m/s. Given that the mass of a proton is 1.672 × 10–27 kg.
ANS: 1.33 × 10–12 m
d
P5.7 Show that cos 2x is not an eigenfunction of the operator dx
but is an
2
eigenfunction of operator dxd 2 . Obtain the eigen value.
ANS: −4
P5.8 Two observables B and C have corresponding operators B̂ and Ĉ. Both
operators have a common set of eigenfunctions. Show that the two operatiors
commute.
P5.9 Show that the eigen values of a hermitian operator are real.
P5.10 Give reasons why ϑ(y, t) = C ye−ky e−iwt is an acceptable quantum

2
mechanical wavefunction while χ (x, t) = C xe−iwt is not a valid quantum

mechanical wavefunction. C, k and w are arbitray constants.
SA5.1 Name two examples where the classical physics failed. State some distin-
guishing features of Schrodinger, Heisenberg and Dirac picture of quantum
mechanics.
SA5.2 How the state of a system at a given instant is specified in quantum
mechanics? What physical significance may be associated with the wave-
function of a system?
SA5.3 Take an example of a suitable wavefunction to explain the essential
characteristics of a valid wavefunction.
SA5.4 Give the statements of the postulates of quantum mechanics.
SA5.5 What is an operator? Define a hermitian operator.
SA5.6 What is the utility of Schrodinger equation? Give properties of this equation.
SA5.7 Write Schrodinger’s time-independent equation and discuss under what
conditions the equation may be used.
SA5.8 Define at least four algebraic operations of operators. When do two
operators commute?
SA5.9 Two hermitian operators Â and B̂ commute, what information about the
eigen values of these operators is conveyed to you?
SA5.10 A system at instant t is defined by the wavefunction ψ(r, t). An operator
Â operating on the system gives three discrete and non-degenerate eigen
values a1 , a2 and a3 for a given dynamic variable that corresponds, respec-
tively, to the eigen states φ1 (r, t), φ2 (r, t) and φ3 (r, t) of the operator. Is it
possible to write ψ(r, t) in terms of φ1 (r, t), φ2 (r, t), and φ3 (r, t)? If yes,
write the expression.
SA5.11 Differentiate between the Eigen value and the expectation value of an oper-
ator. Explain the difference in the outcome of a quantum and a classical
measurement of a dynamic variable.
SA5.12 Distinguish between bound and scattering states of a potential and explain
characteristics of bound states and their wavefunctions.
SA5.13 Give reasons why a plane wave cannot represent a free particle.
SA5.14 A particle is incident on a potential barrier of height greater than the
kinetic energy of the particle, without derivation, write expressions for the
particle wavefunctions in different regions of space and give the boundary
conditions that these wavefunctions must obey.
SA5.15 Write a note on uncertainty principle.
SA5.16 Discuss the principle of correspondences and the logic behind it.
MC5.1 Operators that are their own hermitian conjugate are called
(a) null operator, (b) linear operator, (c) hermitian operator (d) delta
function
ANS: (c)
MC5.2 If for an operator (− F̂) = − F̂ † , then F̂ is
(a) hermitian, (b) anti-hermitian, (c) null operator and (d) unitary operator
ANS: (b)
MC5.3 The probability for obtaining the eigen value (aj ) for a system in quantum
state ψ(r, t) and eigenfunctions ϑ j (r, t) is given by
| +∞ |2
|{ |
( ) | |
P a j = || ϑ j (→ r , t))d3 x ||
r , t)∗ (ψ(→
| |
−∞
Only when
(a) ψ(r, t) is normalised, (b) ϑ j (r, t) are normalised, (c) eigen states are
non-degenerate and (d) operator is hermitian
ANS: (a), (b), (c) and (d)
MC5.4 A particle in quantum mechanics is represented by a wave packet. The
size of the wave packet in one-dimensional space is defined by the width
∓Δx, where the intensity of the packet changes to
1
(a) e2 -times of its central value, (b) e times of central value, (c) e
times of
central value and (d) √1e times of central value
ANS: (d)
MC5.5 The ratio of the ground state energies in a 2D and 3D infinite potential
wells of equal sides is
(a) 3/5 (b) 2/3 (c) 3/2 (d) 5/3
ANS: (b)
MC5.6 Transmission coefficient for a particle of mass m and energy E through a
potential barrier of height V 0 (V 0 > E) and width ‘a’ is given as
( √ )
( ) − 2a 2m (V0 −E )
ℏ
(a) ( 1
a2 V
) (b) ( 1
2
) (c) 16E
V0
1− E
V0
e (d)
1+ 2ℏ20 1+ a2ℏ
2V
[ ( √ ) ]−1
0
( ) −
2a 2m (V0 −E )
ℏ
16E
V0
1− E
V0
e
ANS: (c)
MC5.7 A particular excited state of a system has a width of 1.0 × 10−7 eV. The
mean life of the state will be of the order of
(a) 3.3 × 10−9 s (b) 3.3 × 10−7 s (c) 3.3 × 10−5 s (d) 3.3 × 10−3 s
ANS: (a)
MC5.8 In a classical experiment the speed v of electron was measured three times
to get experimental values as 100 m/s, 110 m/s and 120 m/s. The expec-
tation value of speed ⟨v⟩ obtained using quantum mechanical tools was
found to be ⟨v⟩ = 115 m/s. The expectation value 112 m/s should be
compared with the experimental value
(a) 100 m/s, (b) 110 m/s, (c) 120 m/s and (d) none of the experimental
values
ANS: (b)
MC5.9 If T and R, respectively, denotes the transmission and reflection coeffi-
cients for a potential barrier of height v0 and breadth a for an incident
particle of mass m and energy E ∼ V0 , the ratio R/T is given by
ma 2V0 ℏ2 2ℏ2 ma 2 V0
(a) 2V0 ℏ2
, (b) ma
, (c) ma 2 V0
and (d) 2ℏ2
ANS: (d)
MC5.10 The energy of the first excited state of a particle confined in a 1D infinite
potential well is 12 eV. The fifth excited of the system will be at energy
(a) 27 eV, (b) 48 eV, (c) 70 eV and (d) 108 eV
ANS: (d)
LA5.1 Give reasons why the state function for a free particle cannot be represented
as a plane wave. Explain how a wave packet representation may remove all
problems associated with plane wave representation. How the wave packet
is formed and what are measures for the size of the wave packet in k-space
and spatial space?
LA5.2 Derive expression for the energy levels of a one-dimensional infinite poten-
tial well and show that the eigenfunction for different energy levels shows
nodes depending on the degree of excitation.
LA5.3 A particular experiment of measuring the value of a dynamic variable was
carried out both quantum mechanically and classically. What differences do
you expect in the results obtained in two sets of experiments. What is the
expectation value in quantum experiment and with which value of classical
result you will compare it?
LA5.4 Define a potential barrier and discuss the quantum mechanical tunnelling of
a 1D potential barrier by a particle. Obtain expressions for the wavefunctions
in different region of space around the barrier and apply boundary conditions
to obtain transmission coefficient.
LA5.5 Define (a) the sum, (b) the product of two operators, (c) multiplication of an
operator by a complex number and (d) the commutation of two operators.
Also show that that anticommutator of two hermitian operators is hermitian.
LA5.6 Define
[ ]the commutator of two operators and obtain the value of commutator
Δ
Â, px where the operator Â = x̂ is defined by Âϑ(x) = xϑ(x) and the

Δ Δ
operator px defined as px ϑ(x) = −iℏ dϑ(x)

dx
.
Chapter 6
Quantum Statistics
Objective
Systems having large number of identical particles obey laws of statistics. These
laws assume specific manifestation in case of quantum systems. The distribution
of particles in different quantum energy levels and in different energy states of the
given level and transition from one microstate to the other, etc., are all governed
by quantum distribution laws. Quantum laws of statistics will be introduced in this
chapter. It is expected that a reader will be able to apply laws of quantum statistics
to real cases after going through this chapter.
6.1 Introduction
In Chap. 5 we studied the quantum behaviour of particles; non-interacting particles

confined in a volume of space have discrete non-uniform energy levels, in contrast
to classical physics that predicts continuous energy states/levels.
Quantum statistics or quantum thermodynamics, on the other hand, gives a rational
understanding of statistics in terms of microscopic particles and their interactions.
It deals with systems that contain large number of identical particles or entities, like
a piece of crystalline solid that has large number of identical atoms or molecules
arranged in a specific way or a gas contained in a volume that also has large number
of molecules that are in random motion. Systems of identical entities or particles are
called assemblies, and a system that has large number of subsystems or assemblies
of identical entities is termed as an ensemble.
An assembly has, in general, very large number of entities/particles, for example,
one microgram of an inert gas like argon at some finite temperature and pressure
contains more than 1.5 × 1016 argon atoms moving randomly, colliding with each
other and with the walls of the container. In such a situation it is neither possible nor
required to know the velocity or energy, etc. of each gas molecule at each instant.
What is, however, important to know is the number of gas molecules that have their
https://doi.org/10.1007/978-3-031-32084-2_6
318 6 Quantum Statistics
velocities or energies, etc. in a given range, say, velocities between v and (v + Δv)
or energies between some value ∈ and (∈ + Δ∈) and how this number changes with
time. Conversely, one will like to know as to how a given amount of energy E will
get distributed into different groups of particles in the system. This information is
contained in what is called the distribution function, so there may be several types
of distribution functions like the velocity distribution function or energy distribu-
tion function, etc. for a system. Statistical mechanics or Quantum statistics is
the tool to obtain these distribution functions. Quantum statistics uses the theory
of probability, like the classical statistics, but assumes discrete, i.e. non-continuous
values for physical variables like velocity, energy, momentum, etc. Quantum statis-
tics further assumes that an assembly of identical particles or entities may follow
different kinds of statistics, like Fermi–Dirac, Bose–Einstein or Maxwell–Boltz-
mann statistics. These statistics differ from each other as to how the entities of the
system may be distributed into various energy levels and energy states in a level.
In statistical mechanics, the science of bulk matter is an incomplete and evolving
science. New ideas and concepts permit a fresh approach to old problems. With new
concepts one looks for features ignored in past and expect exiting results. Important
new concepts are: deterministic chaos, fractals, self-organised criticality (SOC),
turbulence and intermittency. These words represent large fields of study, all using
quantum statistics, which have changed how we view nature. Disordered systems,
percolation theory and fractals find applications not only in physics and engineering
but also in economics and other social sciences.
Thus having obtained the required distribution function from the quantum statis-
tics, one analyses the distribution function to obtain the value of a parameter called
the ‘partition function’. Partition function, which depends on the type of the statis-
tics obeyed by the constituent particles of the system, is the most important parameter
from the point of view of quantum statistics. Considerable efforts are put in obtaining
an appropriate partition function for a given system. Partition function, which is like
the heart of quantum thermodynamics, may be used to obtain physical observables
of the system, i.e. the quantities like temperature, pressure, volume, specific heat
capacities, entropy, etc. that may be measured experimentally. In this chapter we will
see how one can use statistical mechanics (or quantum statistics) to get distribution
functions and also how these distribution functions can be further analysed using
the tools of quantum thermodynamics to yield the all important partition functions
which in turn provide values of the required system observables.
6.2 Application of Quantum Statistics (Statistical

Mechanics) to an Assembly of Non-interacting Particles
Statistical mechanics may be applied to solve problems related to real systems that
contain large number of identical entities or particles. The formalism of statistical
physics may be developed both for the classical systems andfor quantum systems.
6.3 Energy Levels, Energy States, Degeneracy and Occupation Number 319
The resulting energy distribution and calculating the values of physical observables
is simpler in the classical case. However, the formulation of the method is more
transparent in the quantum mechanical formalism. In addition, the absolute value
of the entropy without any undermined constant and the behaviour of the entropy
when absolute temperature approaches zero, may be obtained only in the quantum
mechanical treatment. In the following sections we will see how quantum statistics
may be applied to an assembly of non-interacting (or free) particles.
6.3 Energy Levels, Energy States, Degeneracy

and Occupation Number
In Chap. 5 it was shown that a particle confined in an infinite potential of width ‘a’
has discrete set of non-uniform energy levels with energies,
π 2 2
∈ j = n 2j (6.1)
2ma 2
If V denotes the volume of the space in which the particle is confined, then V = a 3
and Eq. (6.1) may be written as
2 (2π )2 − 2
∈ j = n 2j V 3 (6.2)
8m
h
In Eq. (6.2), is rationalised Planck’s constant = 2π = 1.05457 × 10−34 J s.
The integer n j is made up of three independent integers, n x , n y and n z , called the
quantum numbers, such that
n 2j = n 2x + n 2y + n 2Z (6.3)
Each of these n x , n y and n z can have non-zero integer values like 1, 2, 3, …, etc.
So far terms ‘energy level’ and ‘energy state’ have been used interchangeably,
as in most cases energy levels were non-degenerate; however, now onwards energy
level and energy state will have definite meaning as specified here.
The value of n2j defines an energy level of the system. Each energy level may
have one or more energy states. A set of the different values of quantum numbers
n x , n y and n z subject to the condition given by Eq. (6.3) defines the number of
states of the given energy level.
Since the minimum value that n x , n y and n z may have is 1, the minimum value of
n 2j = 12 + 12 + 12 = 3 and, therefore, the lowest energy level has the energy, ∈1 ,
given as,
2 (2π )2 − 2
∈1 = 3 V 3 (6.4)
8m
Table 6.1 Different sets of

nx ny nz
n x , n y and n z that give the
same energy 3 2 1
3 1 2
2 3 1
2 1 3
1 3 2
1 2 3
It may be observed that only one set of n x , n y and n z can give the value 3 to n 2j .
In the language of quantum mechanics it is said that the level ∈1 has only one energy
state. A level that has only one energy state is called a non-degenerate level. The
degeneracy of a level is denoted by g j and is equal to the number of energy states in
the level. The degeneracy g1 of level at energy ∈1 is 1, i.e. g1 = 1.
The next level will be one in which one of the quantum numbers n x , n y or n z has
the value 2 and the other two have values 1. This gives rise to three different sets of
quantum numbers, giving the same value of n 2j = 22 + 12 + 12 = 6. These three sets
are
(n x = 1, n y = 1, n z = 2); (n x = 1, n y = 2, n z = 1) and (n x = 2, n y = 1, n z = 1)
All the three different energy states mentioned above have the same energy ∈2 =
6 (2π )2 − 23
2
8m
V . The level with energy ∈2 has three states and the degeneracy of this
level g2 = 3.
Let us consider the level with energy ∈ = 14 (2π ) 2 2
V − 3 . The six different sets of
2
8m
n x , n y and n z shown in Table 6.1 give the same value of n 2j = 14 and hence the same
energy.
This level, therefore, has sixfold degeneracy or g = 6 for this level. In general the
energy levels are non-equidistant and have different folds of degeneracy.
In particular it may be observed that the three-dimensional energy expression
(6.2) is an equation of sphere of radius R,
2 (2π )2 − 2
∈ j = n 2j V 3
8m
or
8m∈ j 2
n 2j = n 2x + n 2y + n 2z = 2 V 3 = R2
(2π )2
The sphere of radius R is shown in Fig. 6.1. It may be observed in this figure that
the positive non-zero integer values of nx , ny and nz lie in 1/8 quadrant of the sphere.
It means that for large values of energies the density of non-zero positive nx , ny and
nz points essentially fill the volume of the 1/8 quadrant.
Fig. 6.1 Positive non-zero

values of nx , ny and nz lie in
1/8th quadrant of the sphere
One may now treat R or energy ∈ as continuous variable and may obtain the
number of lattice points consistent with energy ≤ ∈ j , which is essentially the volume
of the 1/8 of the sphere. The number of energy states
3/2
1 4 1 2 3/2 1 8m∈ j 2
G(∈) = πR = π R
3
= π 2 V 3 (6.5)
8 3 6 6 (2π )2
And the number of energy states in a thin shell of energy Δ∈ is given by

2 3/2
π 8mV 3 1/2
G(∈, Δ∈) = ∈ j Δ∈ j (6.6)
4 2 (2π )2
In order to have a feel of the magnitude of the degeneracy G(∈, Δ∈), one may
calculate the degeneracy using Eq. (6.6) for molecules/atom moving in a room of size
10 m × 10 m × 10 m at temperature 300 K, taking m ≈ 10−25 kg and Δ∈ j = 0.01∈ j ,
it comes out to be of the order of 1030 .
A quantum mechanical system of N identical non-interacting particles confined
in a given volume of space has many energy levels with each level having a certain
fold of degeneracy. Depending on the properties of the particles, each energy level
accommodates a certain number of particles. If the jth level contains N j particles,
then N j is called the occupation number of the level. The occupation numbers and
the energy of different levels satisfy the following conditions:
N j = N and ∈j Nj = E (6.7)
Fig. 6.2 Schematic representation of energy levels, energy states, occupation no. and degeneracy
of a hypothetical system
Here, E is the total energy of all particles and N their total number.
Figure 6.2 is a schematic representation of the energy levels, energy states, folds
of degeneracy and the occupation numbers of an imaginary assembly. The energy
levels in the figure are represented by horizontal lines, their energies are written on
the left-side vertical scale, and energy states in each level are shown by pink brackets;
the number of energy states in a given level gives the degeneracy of the level. Particles
in different energy states are shown by round dots. Total number of particles in a
level gives the occupation number of the level. It may be observed in the figure that
in some levels there are empty energy states that contain no particle.
SAQ: What creates degeneracy in an energy level?
6.3.1 Distinguishable and Indistinguishable Particles
Properties of particles in a system decide their distribution in different levels and

energy states. One of the important characteristics of a system of identical particles
is whether particles are distinguishable or not. Identical particles may be indistin-
guishable from each other, like the molecules of a gas. Identical molecules of a gas
are in constant motion; therefore, it is not possible to put a mark on one particular
molecule and identify it at all times. Hence, molecules of a gas are indistinguishable
because they are non-localised. On the other hand, in a crystalline solid, atoms or
molecules are identical, but they may be differentiated or distinguished from one
another on the basis of their location in the crystalline lattice. As such, the atoms or
molecules in a crystalline solid are distinguishable as they are localised. In general
non-localised particles are indistinguishable, while localised entities are distinguish-
able because of their fixed location. Atoms/molecules of a paramagnetic salt if put
in an external magnetic field align either parallel or antiparallel to the applied field
and, therefore, may be distinguished from each other through their orientation in
the external field. Similarly, nucleons (neutrons and protons) in a nucleus may be
distinguished from each other on the basis of the orientation of their spins and so
they are distinguishable. The fact that particles of an assembly are indistinguishable
or distinguishable, as we will see, plays an important role in quantum statistics.
6.3.2 Macrostate
A given assembly of identical non-interacting particles has a specific structure of

energy levels and energy states in each level. If the total number of particles N
and total energy E of the system is known it is possible to distribute particles in
different energy levels such that the conditions put by Eq. (6.7) are fulfilled. It is
evident that in general there may be several different ways in which particles may be
distributed in different energy levels and each way of particle distribution satisfies Eq.
(6.7). Each such particle distribution constitutes a Macrostate of the system which
is characterised by a set of occupation numbers for different levels.
Let us explain the concept of the Macrostate by taking an example. Suppose there
is an assembly of non-interacting identical particles that are indistinguishable. For
simplicity we assume that there are only 5 particles, that is N = 5 and that there are
four energy levels available to the assembly, respectively, at energies 0, ∈, 2∈ and 3∈.
Let the total energy E of the system be 10∈. The four possible ways of particle
distribution in different energy levels, consistent with boundary conditions: j N j =
5 and j ε j N j = 10ε are shown in Fig. 6.3.
It may be noted that each of the four Macrostate is characterised by the occupation
numbers of its energy levels. Thus it may be said that ‘The Macrostate of a system
is defined by the number of particles in each energy level of the system’. This
essentially means that if the occupation numbers of all energy levels of a system are
known, the Macrostate of the system is completely defined or known. For example,
the Macrostate (a) in Fig. 6.3 is completely defined by the set of occupation numbers
(N1 = 0, N2 = 0, N3 = 5, N4 = 0) and Macrostate (d) by the set of occupation
numbers (N1 = 0, N2 = 1, N3 = 3, N4 = 1).
It may be observed that while defining Macrostate no consideration is given to the
number of states in each level (i.e. the degeneracy g) and to the number of particles
that may be accommodated in each state.
Fig. 6.3 Particle distribution in four microstates with different sets of occupation numbers
6.3.3 Microstates
As already mentioned, the energy and occupation numbers of different energy levels
as well as the degeneracy (or the number of states associated with a given energy
level), etc. are all decided by the nature of particles in a given assembly. Further, the
number of Macrostates and their configurations for a given assembly are characterised
only by the occupation numbers of different energy levels, and Macrostates do not
depend on the degeneracy of energy levels.
In order to understand the concept of microstates, let us consider the Macrostate
(a) of Fig. 6.3 that has all the five particles in level-3 and is characterised as (N 1 = 0,
N 2 = 0, N 3 = 5 and N 4 = 0). Let us assume that the degeneracy of level-3 is three;
i.e. g3 = 3. Now the five particles in level-3 may be accommodated in three energy
states, shown by coloured brackets in Fig. 6.4, in different ways. If it is assumed
that there is no restriction on the number of particles that may be accommodated
in a state, then six different configurations of particles in level-3 may be shown by
six figures in the lower part of the figure. Each of these configuration is called a
microstate of the Macrostate (N 1 = 0, N 2 = 0, N 3 = 5, N 4 = 0) shown at the top.
It may be noted that the six microstates shown in the figure are not the only
microstates associated with the given Macrostate, there may be many more, and the
actual number will depend on the rules that govern the distribution of particles in
different energy states. A microstate of the system is defined not only by the number
of particles in a level but also by the number of particles in each energy state of
Fig. 6.4 Some microstates corresponding to the macrostate (N 1 = 0, N 2 = 0, N 3 = 5 and N 4 = 0)
each level. Since shifting of a single particle from one energy state of a given level
to another energy state of the same level results in a new microstate, the possible
number of microstates associated with a given Macrostate is very large.
In conclusion it may be said that an assembly of identical non-interacting particles
may have several Macrostates and that each Macrostate may have a very large number
of microstates associated with it.
SAQ: Which physical process may change the microstate of a system?
6.3.4 Time Evolution of an Assembly
Let us assume that a given assembly of non-interacting particles has N-number of

Macrostates designated as M 1 , M 2 , M 3 , … M N . Further, let us assume that Macrostate
M 1 has N 1 microstates designated as k1M1 , k2M1 , k3M1 . . . k NM11 ; Macrostate M 2 has N 2

microstates designated as k1M2 , k2M2 , k3M2 . . . k NM22 ; and so on.
M
Let at a given instant of time, say, t 0 the assembly is in microstate kr j of Macrostate
M j . The assembly will stay in this microstate for a very short time as particles in
same or different states of a given level may collide with each other and will change
M
the distribution of particles shifting to a new microstate k p j of the same Macrostate
M j . As such the system of particles or assembly will quickly shift from one to
the other microstates of the Macrostate M j , staying in each microstate for a very
short time. In the meantime while passing through different microstates of the given
Macrostate, it may so happen that a particle belonging to a given level collides with
a particle belonging to a different level. Since the two colliding particles belong to
two different levels, the collision results in changing the energy of colliding particles
by large amount and particles shift from one level to another level. This results in
the change of the Macrostate, say from M j to M ( j+1) . Now the assembly will stay in
the new Macrostate M ( j+1) for some time and pass through the various microstates
associated with this Macrostate, till collision between particles from two different
levels changes the Macrostate. Since collision between particles of same energy level
is more frequent and involves very small change of particle energies, the assembly
rapidly explores all the microstates associated with the given Macrostate and shifts to
the other Macrostate only when particles from different levels collide to incur large
change in particle energy. This is shown in Fig. 6.5.
Fig. 6.5 Assembly of particles in different macrostates and microstates

6.4 Quantum Statistical Probability of a Macrostate 327
6.3.5 Postulate of Equal a Prior Probability of All Microstates
A fundamental axiom of quantum statistics is that each microstate of an isolated

system is equally likely or has same probability of occurrence. Since this assumption
is made before hand and is for all microstates irrespective of their Macrostate, it is
called the ‘a prior’ assumption. There may be three aspects of the assumption: (i) the
time for which the system lives in a microstate is same for all microstates, irrespective
of to which Macrostate the microstate belongs; (ii) over a given time interval that is
sufficiently large, the system passes through each microstate same number of times;
and (iii) if there are very large number say N, of exactly identical systems and at a
given time N 1 out of them are in some microstate, then the number of systems in
each of the other microstate of the system will also be N 1 .
A system is most stable in its equilibrium state and, therefore, stays in the equilib-
rium state for infinitely long time. It means that the equilibrium state has the largest
number of microstates associated with it.
The assumption of equal a prior probability can neither be derived from some
other more fundamental principle nor be verified by any experiment. The validity of
the principle is based on the correctness of the results derived from it.
In the preceding example of M 1 , M 2 , M 3 , …M N Macrostates with Macrostate
M 1 having N 1 microstates, Macrostate M 2 having N 2 microstates and so on, if it is
assumed that the assembly stays in any microstate for a time Δt, then the system will
stay in Macrostate M 1 for time TM1 = N1 Δt, in Macrostate M 2 for time TM2 = N2 Δt,
in Macrostate M 3 for the time TM3 = N3 Δt and so on. The system stays for a longer
time in those Macrostates which have larger number of microstates. Equilibrium state
of a system, where the system stays for infinitely long time, has the largest number
of microstates associated with it.
SAQ: Which physical process may result in the change of a Macrostate?
6.4 Quantum Statistical Probability of a Macrostate
Each Macrostate of a system is assigned a thermodynamic probability or statis-

tical count. The number of equally likely microstates associated with a given kth
Macrostate is called the thermodynamic probability or statistical count of the kth
Macrostate and is denoted by Wk . It may be noted that in quantum thermodynamics
probability of a Macrostate is a number, while generally probability is a ratio or
fraction.
The thermodynamic probability of the system as a whole is equal to the sum of
all microstates and is denoted by Ω,
Ω= Wk
k
6.4.1 System Properties and Average Occupation Number
As shown in Fig. 6.4, the microscopic structure of the system changes almost
continuously with time as the system moves through different Macrostates. Since
Macrostates live for short times, it is not possible to measure the physically important
system properties in a particular Macrostate. Actually whenever some measurement
is done the system passes through a large number of Macrostates during the process
of measurement. Therefore, the measured physical quantity is the average value over
a large number of Macrostates of the system. A change in the Macrostate means
change in the occupation numbers of the levels. The measured value of the physical
quantity, therefore, depends on the average values of occupation numbers of different
energy levels of the system. As such values of occupation numbers for different levels
averaged over a large number of Macrostates are of paramount importance from the
point of experimental determination of system properties.
Quantum statistics provides a method to compute the average occupation number
of some level, say level j, of the Macrostate k, which is denoted by N jk . It is quite
obvious that the value of N jk will depend on the probability Wk of the Macrostate k.
The value of Wk depends on the nature of the particles of the system. Quantum
mechanics classifies particles according to the statistical distribution law that is
followed by a group of large number of identical particles. There are three different
statistical distribution laws, namely Bose–Einstein, Fermi–Dirac and Maxwell–
Boltzmann distribution law, one of which is followed by a group of large number
of identical particles. Each distribution law gives a different value of Wk . We first
calculate the value of Wk for systems that follow these three distribution laws.
6.5 The Bose–Einstein Statistical Distribution
Bose–Einstein distribution law is applicable when:

1. Particles are indistinguishable.
2. Any number of particles may be accommodated in a given energy state.
3. States are distinguishable.
Let us consider the jth level of the system and assume that there are gj states in this
level in which N j particles are distributed. Since there is no restriction on the number
of particles in a state, there may be a large number of ways in which N j particles may
be distributed in gj states. We will now calculate the number of these possible ways
of distribution of particles. In order to make these calculations we designate particles
by lower case letters a, b, c, …. As a matter of fact this representation of identical
and indistinguishable particles by different distinguishable letters is wrong, but it is
being done to make calculations simple, further, as you will see the indistinguishable
nature of particles will be maintained in calculations. The different states in level j
6.5 The Bose–Einstein Statistical Distribution 329
are represented by numbers 1, 2, … gj as shown in Fig. 6.6. One possible distribution

of particles in level j may be the one shown in Fig. 6.6, which may be written as,

{(1)ab}{(2)}{(3)cde} . . . g j lm (6.8)
In Eq. (6.8) curly brackets represent states and lower case letters the particles.
The sequence shown in Eq. (6.8) is made up of gj numbers (representing states)
and N j particles that mean a total number of elements in the sequence are (gj +
N j ). Any sequence of these (gj + N j ) elements of the type indicated in Eq. (6.8)
gives a way of distribution of particles in different states of the level. But there is
one condition on the valid sequence that represents particle distribution is that the
sequence must start with a number representing the state. A sequence of the type
{ab(2)}{(1) f gl} . . . g j lm is not valid as it starts with a letter and not numbers.
If the sequence starts with one out of the g j numbers, the number of remaining
elements becomes g j + N j − 1 . The number of different ways in which these
remaining
elements
g j + N j − 1 may be arranged is factorial gj + Nj − 1 ,
i.e. g j + N j − 1 !. To compute the total numberof valid sequences
it is required
to multiply thenumber of different
ways in which g j + N j − 1 elements can be
arranged (i.e. g j + N j − 1 !) by g j , any one of which may be the first element
of the valid sequence. Thus the total number of different ways in which N j particles
may be distributed in g j states is,

N total = g j N j + g j − 1 ! (6.9)
The number N total contains

sequences in which arrangements like the following
{(1)ab}{(2)}{(3)cde} . . . g j lm and {(1)ba}{(2)}{(3)cde} . . . g j lm , etc. are
treated as different sequences. However, the particles are indistinguishable; hence
such sequences are identical and do not give a new way of particle distribution. We
Fig. 6.6 Distribution of N j identical particles in gj distinguishable states

thus observe that in N total some; otherwise identical sequences have been counted as
different sequences. This has happened because we assigned distinguishable letters
a, b, c, etc. to the undistinguishable particles. It is, therefore, required to correct N total
for this over counting. The number of particles is N j , and the number of different
ways in which these particles can be arranged is N j!. Hence, correction for this over
counting may be applied by dividing N total by N j !.
Similarly, sequences like
{(1)ab}{(2)}{(3)cde} . . . g j lm and
{(2)}{(3)cde} . . . {(1)ab} . . . g j lm have also been counted as different sequences
in N total . But actually these are not two different sequences. It may be noted that
states are distinguishable, which means that number 1, 2, 3, … are different but
at which location in the sequence a particular number occurs is not important. As
shown in the two sequences above the state {(1)ab} appears in the first location of
the sequence or it appears at any other location does not matter so long the number
of particles in the state remains same. As such, the two sequences shown above
refer to the same distribution. There are g j numbers, and the possible ways in which
they may be arranged are g j !. Correction for this over counting may be applied by
dividing N total by g j !
Finally, the corrected number of different ways in which N j indistinguishable
particles may be distributed in g j states or the number of different distributions for
the jth level ω j is,

N total gj gj + Nj − 1 ! gj + Nj − 1 !
ωj = total
Ncorrected = = = (6.10)
(g j !) N j ! (g j !) N j ! gj − 1 ! Nj!
Before proceeding further let us check the correctness of Eq. (6.10). For simplicity
we assume that there are only three particles (N j = 3) and only 3 states (g j =
3); then according to Eq. (6.10), the number of different ways in which these 3
indistinguishable particles may be distributed in 3 states is,
5!
ω(3, 3) = = 10
(2!)(3!)
These ten different ways of particle distribution are shown in Fig. 6.7, where dots
represent particles.
Application of formula given by Eq. (6.10) for calculating the number of ways
of particle distribution to the case of a non-degenerate level needs a mention. For
a non-degenerate level g j = 1 and, therefore, in the denominator of expression
(g j +N j −1)! (1+N j −1)!
{(g j −1)!}( N j !) one gets {0!}( N j !) = 0! . Now for a non-degenerate level there is only
1
one way of distributing indistinguishable particles, that is all particles are in the same
state. Hence, in order to match formula given by Eq. (6.10) with the experimental
fact we should use the convention that 0! = 1. This convention will also make the
formula valid for the state which is empty and has no particle. For an empty state N j
= 0 and wempty = (g −1j !(0!)
0+g −1)!
= 0!1 = 1.
( j )
6.5 The Bose–Einstein Statistical Distribution 331
Fig. 6.7 Ten different ways of distributing three particles in three different states
Equation (6.10) gives the number of possible ways in which particles may be
distributed in any level. Suppose in one of the levels particles are distributed in one
of the ways given by Eq. (6.10). Then in each of the remaining levels particles may
be distributed according to any one of the distribution out of those specified by Eq.
(6.10). Therefore, the total number of possible distributions or the statistical (ther-
modynamic) probability of any Macrostate in Bose–Einstein distribution is given
by,

gj + Nj − 1 !
W Bose-Ein
= (6.11)
j
gj − 1 ! Nj!

The symbol j f ( j ) means the forming of products of each term of function f ( j)
for each value of j.
Suppose there is a system of three energy levels, such that in level-1 there are 2
particles and 2 states, in level-2: 3 particles and 3 states and in level-3: 1 particle
and 3 states. In other words the occupation numbers and the folds of degeneracy of
level-1, level-2 and level-3 are, respectively, N 1 = 2, g1 = 2; N 2 = 3, g2 = 3; and
N 3 = 1, g3 = 3. Once the occupation numbers of the levels have been fixed it means
that the Macrostate is fixed. The thermodynamic probability of this Macrostate for
Bose–Einstein distribution can be calculated using formula given by Eq. (6.11) as

follows:

gj + N − 1 !
W Bose-Ein =
j
g j − 1 ! (N !)

(g1 + N1 − 1)! (g2 + N2 − 1)! (g3 + N3 − 1)!
=
{(g1 − 1)!}(N1 !) {(g2 − 1)!}(N2 !) {(g3 − 1)!}(N3 !)

3! 5! 3!
= = {3}{10}{3} = 90
(1!)(2!) (2!)(3!) (2!)(1!)
6.6 The Fermi–Dirac Statistical Distribution
Indistinguishable particles that obey Fermi–Dirac statistics follow this distribution.

In this distribution it is assumed that:
(1) Particles are indistinguishable.
(2) Not more than one particle can be accommodated in one state. This means that
the number of particles N j in jth level cannot be larger than the number g j of
states in level j.
(3) States are distinguishable.
Let there be N j particles in level j which are distributed in g j states, where
g j ≥ N j . Again we assign distinguishable numbers 1, 2, 3, … to distinguishable states
in a level and lower case letters a, b, c… to indistinguishable particles. Correction for
choosing distinguishable letters for indistinguishable particles will be applied later
on. A possible sequence representing the arrangement of particles in the level j may
be

1()2(a)3(c)4() . . . g j (m)
The distribution shown above has no particle in first state, a particle in state-2, a
particle in state-3, no particle in state-4… and one particle in the last state g j .
The problem of particle distribution may be looked in the following way:
Suppose initially all the g j states in level j are empty. We now take one particle
(out of the total N j particles) and put it in one of the states. This first particle may be
put in any of the gj states. That means that for the first particle there are gj different
ways of filling the states. If n i denotes the number of ways in which the ith particle
can be filled, then n1 = g j .After putting the first particle in any state, the number
of
empty states left is g j − 1 . Second particle may now be put in one of the g j − 1
states. It means that the number of different
ways
n 2 in which second particle can
be filled in remaining states is n 2 = g j − 1 . Continuing the same argument, the
number of ways in which the third, the fourth and so on up to nth particle filling will
6.6 The Fermi–Dirac Statistical Distribution 333

be, respectively, n 3 = g j − 2 , n 4 = g j − 3 , . . . n n = g j −
n + 1 . Thenumber
of ways in which the last N j th particle can be filled is n N j = g j − N j + 1 .
The total number of ways of distributing particles in all states n total = n 1 × n 2 ×
n3 × . . . n N j .
Or

n total = g j × g j − 1 × g j − 2 × . . . g j − N j + 1
gj!
n total = (6.12)
gj − Nj !
Corrections for over counting of ways of distribution

(i) The particles are indistinguishable, and therefore, sequences like [1(a) 2()
3(b)…] and [1(b) 2() 3(a)…] and [1(d) 2() 3( f )…}, etc. that have been counted
as different sequences in n total basically refer to only one way of distribution.
Correction for this over counting may be applied by dividing n total by N j ! which
is the number of different combinations of a, b, c, … N j .
It may be noted that in present counting of different ways we have not counted
[1(a) 2() 3(b)…] and [2() 1(a) 3(b)…], etc. as different sequences. What has been
done is to fix the locations of all energy states and fill them with particles one
after another. Hence no correction for this needs to be applied.
Therefore,
The number of ways in which N j indistinguishable particles obeying Fermi–Dirac
distribution may be distributed in g j states in level j is,
n total gj!
ωj = = (6.13)
Nj! Nj! gj − Nj !
Once again, to test the correctness of Eq. (6.13) we calculate the number of
different ways in which three particles can be distributed in three energy states of
level j, when particles follow Fermi–Dirac statistics. In this case, g j = 3 and N j = 3.
Substituting these values in Eq. (6.13) one gets,
3!
ωj = = 1.
3!0!
It may be noticed that if particles obey Bose–Einstein statistics then three particles
may be distributed in ten different ways in three energy levels. On the other hand
they may have only one way of distribution if they obey Fermi–Dirac statistics.
Now for any one of the ω j arrangement of particles in a given level, there are
ω j ways of distribution of particles in any other level. Therefore, the number of
ways in which fixed number of particles in each level may be distributed in different
energy states of each level, that is the thermodynamic probability of a Macrostate in
Fermi–Dirac statistics is,
gj!
W Fermi-Dirac = (6.14)
j
Nj! gj − Nj !
SAQ: What are the main points of difference between Bose–Einstein and Fermi–
Dirac statistics?
6.7 The Maxwell–Boltzmann Statistical Distribution
Particles that obey Maxwell–Boltzmann statistics are

(1) Distinguishable.
(2) Any number of particles can be accommodated in an energy state.
(3) Energy states are also distinguishable.
When particles are distinguishable, the number of different ways of distributing
particles in different levels and in states of the same level becomes very large. It is
because the same number of particles in a given state may be put in many different
ways if particles are distinguishable, while there is only one way of putting a given
number of particles in a given state if particles are indistinguishable. For example,
if in state-1, which we denote by number 1, there are three particles that are indis-
tinguishable then sequences 1(abc), 1(bcd), 1(mna), 1(cba), etc. are all equivalent
because there is no way to make a distinction between a, b, c, … m, n, etc. However,
if particles are distinguishable, which means that a, b, c, d, m, n… are all different
then first three sequences are different, and each of them corresponds to a new way
of particle distribution. However, sequences 1(abc) and 1(cba) are equivalent both
when particles are distinguishable or not because ordering of particles in a state does
not create a new microstate.
Another reason for the large number of ways in which distinguishable particles
may be distributed in different levels and states of a level is the fact that if a particle
say ‘m’ moves from an energy level j to another level k and a particle ‘n’ from
level k goes to level j, a new microstate is created even if the number of particles
(occupation numbers) in levels j and k remains same. On the other hand no new
microstate is created if the occupation numbers of levels do not change when particles
are indistinguishable.
Suppose there are N distinguishable particles in all, and they are to be distributed
say, in levels j1 , j2 , j3 , … such that level j1 has N 1 particles; level j2 , N 2 particles;
level j3 , N 3 particles and so on. The total number of particles in all levels is N, so that
N 1 + N 2 + N 3 + · · · = N. We now calculate the number of different ways in which
this distribution can be done. We mark N locations and try to place N distinguishable
particles in these locations one after another and count in how many different ways
this can be done. Suppose we take a particular particle and place it in one of the N
locations. The first particle can be placed in N different ways in these N locations.
The second particle will have now (N − 1) locations and can be placed in one of these
available locations in (N − 1) different ways. The third particle will have only (N −
6.7 The Maxwell–Boltzmann Statistical Distribution 335
2) vacant locations and can be placed in one of these locations in (N − 2) different

ways. The 4th, the 5th and so on successive particles may be put, respectively, in
(N − 3), (N − 4), (N − 5)… different ways. The last particle will have only one
vacant place and can be but only in 1 way. Therefore, the total number of different
ways in which these N particles may be placed in N locations is N (N − 1) (N −
2) (N − 3)…1 = N!. Since each arrangement of N particles makes a sequence, the
number of different sequences of N particles is N!.
Next suppose that out of these N particles placed in some sequence, the first N 1
particles go to fill states in level j1 , next N 2 particles in level j2 , next N 3 particles in
level j3 and so on. We have calculated that the number of different possible sequences
of N different particles is N!, and if each sequence is broken down into groups of N 1 ,
N 2 , N 3 ,…. particles, the number of different sequences for each group of particles
will also be N!. Now in N! different sequences for each group of particles, some
sequences will be those in which same particles will be placed at different locations.
For example, if N 1 = 5, then sequences of the type (i) a c d b e, (ii) c a d b e, (iii) b
a d e c …, though counted as different sequences in the number N!, are not different
as ordering of distinguishable particles does not make a new sequence. Therefore,
for each group of N 1 , N 2 , N 3 , … particles over or excessive counting of different
sequences have occurred. In order to apply correction for this over counting, let
us calculate the number of different sequences that can be made by putting same
particles at different positions. For example, in the case N 1 = 5, letters a b c d e may
be arranged in 5! = 120 different sequences, only three of which are shown above.
If the group of particles consists of N 1 particles, the number of different sequences
of same particles will be N 1 !. Similarly group of N 2 particles will have N 2 ! identical
sequences and so on. Correction for excess counting may, therefore, be applied by
dividing N! by (N 1 ! × N 2 ! × N 3 ! × · · · ).
Hence, the number of different ways in which N 1 , N 2 , N 3 … particles out of total
N distinguishable particles may be filled, respectively, in levels j1 , j2 , j3 , … without
repeating the sequence of same particles is given by,
N! N!
N level = = (6.15)
[(N1 !)(N2 !)(N3 !) . . .] j Nj!
Next we calculate the number of different ways in which N j distinguishable parti-

cles in level-1 may be distributed in g1 energy states. If we assume that each energy
state is like a box, then we have g1 distinguishable or different boxes and N 1 distin-
guishable particles. At random we take one particle and put it in one of the box, this
may be done in g1 different ways. Next we take any other particle, and this particle
can also be put in box-1, or box-2 or in any one of the g1 boxes. Since there is no
restriction on the number of particles in any state, including the state in which first
particle has already been put; the second particle also has g1 different ways of filling
the energy states. As a matter of fact each of N 1 particles has g1 different ways of
filling the energy states. Therefore, the total number of ways in which N 1 particles
may be filled in g1 states in level-1 is given by,
N l−1 = (g1 ) × (g1 ) × (g1 ) × . . . N1 terms = g1N1
Similarly, the number of different ways in which N 2 particles in level-2 may be

distributed in g2 states is N l−2 = g2N2 and so on. Therefore, the total number of
different ways in which particles may be distributed in different states of different
levels, N state , is given by,
Nj
N state = N l−1 × N l−2 × N l−3 . . . = g1N1 × g2N2 × g3N3 . . . = gj (6.16)
j
The thermodynamic probability W Maxwell-Boltzmann (which is equal to the total

number of microstates) of a Macrostate in Maxwell–Boltzmann distribution may be
obtained by multiplying N level with N state .
⎛ ⎞
N! g Nj j
W Maxwell-Boltzmann = ⎝ gj j⎠ = N!
N
(6.17)
j Nj! j j
Nj!
As an example, let us calculate the thermodynamic probability of a Macrostate for

a system that obeys Maxwell–Boltzmann statistics and in which five distinguishable
particles are distributed in two levels with N 1 = 3, g1 = 3 and N 2 = 2, g2 = 4.
Using Eq. 6.17 one gets,

33 × 24
W Maxwell-Boltzmann
= 5! = 1440
3! × 2!
6.8 Relation Between Entropy and Thermodynamic

Probability
Suppose there are two independent systems A and B with entropies S A and S B . Let
the thermodynamic probabilities of the two systems be Ω A and Ω B . It is known that
entropies are additive; therefore, the total entropy S total of the two systems put together
is
Stotal = S A + S B (6.18)
Now thermodynamic probabilities Ω A and Ω B of the two systems give respec-

tively the total number of microstates of systems A and B. Since for each microstate
of system A, there will be Ω B number of microstates of system B, and the total number
of microstates Ωtotal of the two system put together will be
6.8 Relation Between Entropy and Thermodynamic Probability 337
Ωtotal = Ω A .Ω B (6.19)
It may be noticed that while entropies are additive, thermodynamic probabilities

are multiplicative. As such, there cannot be one-to-one correspondence between
entropy S and thermodynamic probability Ω. That means,
S A /= constant × Ω A and S B /= constant × Ω B
However, it is possible that entropy S of a system is some function of

thermodynamic probability Ω of the system.
Let us assume that
S = f (Ω) (6.20)
where, f represents some function. Our task is to explore the nature of function f .
It follows from Eqs. 6.18 and 6.19 that
f (Ω A ) + f (Ω B ) = f (Ω A Ω B ) (6.21)
We differentiate Eq. (6.21) with respect to Ω A to get
d f (Ω A ) d f (Ω A Ω B )
+0=
dΩ A dΩ A
or
d f (Ω A ) d f (Ω A Ω B ) d(Ω A Ω B ) d f (Ω A Ω B )
= = ΩB
dΩ A d(Ω A Ω B ) dΩ A d(Ω A Ω B )
or
d f (Ω A ) d f (Ω A Ω B )
= ΩB (6.22)
dΩ A d(Ω A Ω B )
Similarly, when Eq. 6.21 is differentiated with respect to Ω B , one gets
d f (Ω B ) d f (Ω A Ω B )
= ΩA (6.23)
dΩ B d(Ω A Ω B )
When Eqs. (6.22) and (6.23) are multiplied, respectively, by Ω A and Ω B one gets,
d f (Ω A ) d f (Ω B )
ΩA = ΩB (6.24)
dΩ A dΩ B
The two sides of Eq. (6.24) contain functions of two independent variables, and
hence this equation will be true only when the two sides are equal to some constant.
Let this constant be denoted by k B . So that
d f (Ω A ) d f (Ω B ) d f (Ω)
ΩA = ΩB = ··· = Ω = kB
dΩ A dΩ B dΩ
or
dΩ
d f (Ω) = k B
Ω
or
f (Ω) = k B ln Ω
or
S(Ω) = k B ln Ω (6.25)
The numerical value of constant k B , called Boltzmann constant, has been deter-
mined by actually matching the value of entropy with ln Ω and has been found to
be,
R(Gas constant)
kB = = 1.38062 × 10−23 J K−1 (6.26)
A(avogadro’s number)
SAQ: What is entropy? What property of the system is measured by the entropy of
the system?
6.9 The Distribution Function
It has already been discussed that the average occupation numbers play an important
role so far as the system observables are concerned. Since even very small physical
systems at room temperature contain very large number of particles and available
energy levels, it is almost impossible to calculate average occupation numbers by
counting levels and calculating possible ways of particle distributions. The average
occupation numbers are, therefore, determined theoretically, using the distribution
laws of quantum statistics, and expressions for entropy change are borrowed from
classical thermodynamics. The expression for average occupation number per energy
state W g
is called the distribution function. It is possible to derive distribution
functions for different statistical distributions, but that is beyond the scope of the
present discussion. Significance of distribution functions lies in the fact that they
may be directly related to the observables, like temperature, pressure, volume, etc.
of s system in equilibrium.
Distribution functions for some important distributions are given here.
Distribution function for Bose–Einstein distribution
6.9 The Distribution Function 339
NJ 1
= μ−∈ j (6.27)
gj −
e T kB − 1
Distribution function for Fermi–Dirac distribution
NJ 1
= (6.28)
gj (∈ j −μ)
1 + e kB T
Distribution function for classical distribution
NJ 1 (μ−∈ j )
= ∈ −μ = e k B T (6.29)
gj (j )
e kB T
Equation (6.29) tells that the average number of particles per state in every level
exponentially decreases with the energy ∈ j of the level, for a system of indistin-
guishable particles obeying classical statistics. Further, the decrease of the number
of particles per state is faster at low temperature. The same result will be found in
case of the systems that follow Maxwell–Boltzmann distribution, though particles
are distinguishable in this distribution.
Distribution function for Maxwell–Boltzmann distribution
Or

Nm μ−∈m
= Ne T kB
gm
Above equation holds for any level of the system, and generalising it one may
write
μ−∈
NJ j
= Ne T kB
(6.30)
gj
Equation (6.30) gives the desired distribution function.

It is clear from Eqs. (6.29) and (6.30) that distribution functions for classical
distribution and Maxwell–Boltzmann distribution are identical if written in terms of
the partition function.
SAQ: What is the significance of a distribution function?
Solved Examples
SE6.1 An assembly of 5 indistinguishable non-interacting particles has 4-

equidistant discrete energy levels with energies 0, ε, 2ε and 3ε. Calculate
the number of Macrostates of the system if the total energy of all particles
is 10ε.
Solution Macrostates of an assembly are defined by the different sets of possible
occupation numbers subject to the conditions: i Ni = N (no. of particles) =
5 and i Ni ∈i = (Total energyE) = 10∈
These conditions restrict the number of Macrostate to 5 as shown below.
SE6.2 Assuming that level-4, level-3, level-2 and level-1 in the above problem
SE6.1 have, respectively, 4, 5, 2, and 1-folds of degeneracy, and the particles
obey Fermi–Dirac statistics calculate the number of microstates associated
with each of the five Macrostates. In which Macrostate the system will stay
for the longest time?
Solution It is given that g4 = 4, g3 = 5, g2 = 2 and g1 = 1. Also, it is given
that the particles obey Fermi–Dirac statistics. The number of microstates in case of
Fermi–Dirac distribution is given by the relation,
g !
W Fermi-Dirac = j { N !}(gj −N ! where symbols have their usual meaning.
j j j)
Let us calculate the microstates of Macrostate (N 1 = 1, N 2 = 1, N 3 = 0 and N 4
= 3)
1! 2! 4!
W Fermi-Dirac = × × =8
{1!}(0)! {1!}(1!) {3!}(4 − 3)!
Microstates of Macrostate (N 1 = 1, N 2 = 0, N 3 = 2 and N 4 = 2) = 1!

{1!}(0)! ×
{2!}(5−2)! × {2!}(2)! = 30.
5! 4!
Microstates of Macrostate (N 1 = 0, N 2 = 2, N 3 = 1 and N 4 = 2) = {2}!(0)!2!

×
{1!}(5−1)! × {2!}(4−2)! = 30.
5! 4!
Microstates of Macrostate (N 1 = 0, N 2 = 1, N 3 = 3 and N 4 = 1) = {1!}(2−1)!

2!
×
5!
{3!}(5−3)! × 4!
{1!}(4−1)! = 80.
Microstates of Macrostate (N 1 = 0, N 2 = 0, N 3 = 5 and N 4 = 0)= {5!}(5−5)!
5!
= 1.
Since Macrostate (N 1 = 0, N 2 = 1, N 3 = 3 and N 4 = 1) has largest number of
microstates 80, therefore, system will stay in this Macrostate for the longest time.
SE6.3 Every other detail remaining same in problem SE6.2 calculates the
microstates corresponding to different Macrostates if the particles obey
Maxwell–Boltzmann distribution.
Solution It is given that all other details are same that means that there are six
Macrostates, and degeneracy of different levels starting from 1 to 4 are g1 = 1, g2 =
2, g3 = 5 and g4 = 4.
Formula for the microstates in case of Maxwell–Boltzmann distribution is given
as;
g Nj j
W Maxwell-Boltzmann
= N!
j
Nj!
Microstates of Macrostate (N 1 = 1, N 2 = 1, N 3 = 0 and N 4 = 3) =

11 21 43
5! 1! × 1! × 3! = 2560.
Microstates of
Macrostate (N 1 = 1, N 2 = 0, N 3 = 2 and N 4 = 2) =
11 52 42
5! 1! × 2! × 2! = 12,000.
22 51 42
5! 2! × 1! × 2! ] = 9600.
21 53 41
5! 1! × 3! × 1! = 20,000.
5
Microstates of Macrostate (N 1 = 0, N 2 = 0, N 3 = 5 and N 4 = 0)= 5! 5
5!
= 3125.
Problems
P6.1 A system of identical non-interacting particles has 5 particles with total energy
5E, distributed over three energy level of energies 0, E and 2E. Show that the
system may have three Macrostates and designate these Macrostates in terms
of their occupation numbers. Assuming that each level is threefold degen-
erate, calculate the number of microstates associated with the Macrostate
with occupation numbers [N (E=0) = 2, N(E=E) = 1, N (E=2E) = 2} when parti-
cles obey (A) Bose–Einstein statistics, (B) Fermi–Dirac statistics and (C)
Maxwell–Boltzmann statistics.
ANS: (A) 108 (B) 27 (C) 7290
P6.2 In a system of total energy 12∈, five indistinguishable particles obeying

Bose–Einstein statistics are distributed in five equidistant levels of ener-
gies 0, ∈, 2∈, 3∈, 4∈, and all levels have four energy states. Calculate the
number of Macrostates, and microstates corresponding to each Macrostate.
Also calculate the mean occupation number for each level.
ANS: Result is summarised in Table 6.2. There will be 9 Macrostates as
specified in the table.
P6.3 ShowN that both Fermi–Dirac and Bose–Einstein distributions reduce to
gj j
j N j ! when g j >> N j .
SA6.1 What is quantum statistics and what is its significance?

SA6.2 What do you understand by the terms: energy levels, energy states,
degeneracy and occupation number? State the postulate of equal a prior
probability.
SA6.3 Distinguish between the Macro- and microstates of a system. How do they
differ from each other?
SA6.4 What is the significance of average occupation number and how they are
related to with physical observables?
SA6.5 What are the differences between the characteristics of Fermi–Dirac and
Bose–Einstein statistics? Write expressions for number of microstates asso-
ciated with a given Macrostate in cases of Bose–Einstein and Fermi–Dirac
distributions.
SA6.6 What are the characteristics of classical distribution? Write expression for
its distribution function.
SA6.7 Give the relationship between the quantum statistics probability and the
entropy of a system.
SA6.8 Comment on the statement ‘entropy is a measure of disorder of the system’.
SA6.9 Does the nature of energy levels of a system depend on the volume of the
system? What happens to the energy levels of a system when the volume
of the system is increased?
SA6.10 How degeneracy evolves in a system of identical non-interacting particles?
SA6.11 In what respect classical distribution differs from Maxwell–Boltzmann
distribution?
SA6.12 Explain the statement ‘Equilibrium state of a system has largest number of
microstates’.
SA6.13 Why transition from one microstate to another microstate of the same
Macrostate is faster as compared to the transition of the system from one
Macrostate to another Macrostate?
6.9 The Distribution Function
Table 6.2 Particle distribution in different energy levels of each macrostate

Energy of the level 1 2 3 4 5 6 7 8 9 Mean occupation number
↓
4∈ ppp pp pp P p P 1.253
3∈ p Pp pp pppp ppp pp 1.441
2∈ pp P Pppp p ppp 0.941
∈ p pp p 0.571
0 pp P p P p 0.720
Number of microstates → 200 400 640 640 400 140 140 320 200 Total number of microstates = 3080
Note that the sum of mean occupation numbers: 1.253 + 1.441 + 0.941 + 0.571 + 0.720 = 5.0 which is the total number of particles in the system
343
SA6.14 What does the distribution function for a statistical distribution tells? What
is its physical significance?

native. All correct alternatives must be marked for the complete answer/full marks
in such cases.
MC6.1 Sum of average occupation numbers over all energy levels of a system is
equal to
(a) Total number of levels (b) total number of particles (c) average value
of the degeneracy of levels (d) average value of quantum thermodynamic
probability
ANS: (b)
MC6.2 The entropy of an assemble and its quantum thermodynamic probability
have functional relationship that is
(a) Linear (b) binomial (c) exponential (d) logarithmic
ANS: (d)
MC6.3 Experimental values of system observables of an assembly depend
strongly on
(a) Occupation number of the level with highest energy (b) occupation
number of the level with least energy (c) average occupation numbers of
all levels (d) temperature of the assembly
ANS: (c)
MC6.4 A system has four particles which are distributed in three energy levels of
energies, 0, E and 2E. If the total energy of the system is 5E, the number
of Macrostates of the system is
(a) 4 (b) 3 (c) 2 (d) 1
ANS: (c)
MC6.5 In question (MC6.4), if each level has threefold degeneracy and the parti-
cles obey Maxwell–Boltzmann statistics then the maximum number of
microstates associated with a Macrostate will be
(a) 100 (b) 238 (c) 324 (d) 972
ANS: (d)
MC6.6 Which of the following expressions represent the quantum probability of
an assembly of particles obeying Maxwell–Boltzmann distribution law?
Symbols have their usual meaning.
(g j +N j −1)! (b) gj! Nj

gj
(a) j {( j ) }( j )
g −1 ! N ! j { j }( j j )
N ! g −N !
(c)N ! j Nj!
(d) Nr
gr
=
1
μ−∈r
−
e T kB
−1
ANS: (c)
MC6.7 A typical Macrostate of a system of four identical particles has two parti-
cles each in two levels of twofold degeneracy each. The quantum thermo-
dynamic probability of the Macrostate is 1 when particles follow statistical
distribution law A and 96 when B. Statistical distribution laws A and B
are respectively
(a) Bose–Einstein, Maxwell–Boltzmann (b) Fermi–Dirac, Maxwell–
Boltzmann (c) Maxwell– Boltzmann, Bose–Einstein (d) Fermi–Dirac,
Bose–Einstein
ANS: (b)
MC6.8 Which of the following expression represent Fermi–Dirac distribution
function?
(g j +N j −1)! gj! g Nj j
(a) j {(g j −1)!}( N j !) (b) j { N j !}(g j −N j )! (c)N ! Nr
j N j ! (d) gr =
1
μ−∈r
−
e T kB
−1
ANS: (b)
MC6.9 Four identical particles equally distributed in two twofold degenerate
levels have 96 microstates. Nature of the particles and the statistics they
follow are
(a) Distinguishable, Bose–Einstein statistics (b) indistinguishable, Fermi–
Dirac statistics (c) distinguishable, Maxwell–Boltzmann statistics (d)
indistinguishable, classical statistics
ANS: (c)
MC6.10 What happens to a level of energy ∈ of a system when the volume of the
system is doubled?
(a) Shifts to a higher energy by 0.37∈

(b) No shift in the energy of the level
(c) Shifts to a lower energy by 0.74∈
(d) Shifts to a lower energy by 0.37∈
ANS: (d)
LA6.1 What is quantum thermodynamic probability for a system of identical

particles? Establish a relationship between entropy and thermodynamic
probability.
LA6.2 What is meant by the statistical distribution of identical non-interacting parti-
cles in energy levels? How do energy levels and degeneracy of levels develop
in quantum statistics? Name three types of quantum statistics that may be
followed by particles and derive an expression for the number of microstates
corresponding to a given Macrostate for one of these distributions.
LA6.3 Distinguish between (a) energy level and energy state and (b) Macrostate
and microstate. State the law of a prior equal probability and discuss its
different aspects. What event may cause transition from one microstate to
another microstate?
LA6.4 Spell out the characteristics of Maxwell–Boltzmann distribution; derive
expressions for (a) quantum thermodynamic probability; and write the
expression for its distribution function.
Chapter 7
Optical Fiber Communication
Objective
Basics of optical fiber communication for terrestrial transfer of information are
discussed in this chapter. It is expected that after reading this chapter the reader
will be able to understand why optical communication is better and faster than both
the wireless and metallic cable transmissions. He will also appreciate the technique
of optical fiber transmission and its requirements.
7.1 Introduction
The most common method of transferring information from one place to another
place is either using wireless transmission or transmission using metal core cables.
In both these methods, though information in analogue form as well in digital signal
form may be transferred but digital mode is preferred on account of its transmis-
sion reliability and less affect from surrounding environmental conditions. In digital
transmission, the required information is first converted into digital electronic pulses,
mostly voltage pulses, which are then transmitted through metal core cables for terres-
trial transmission or are made to modulate a high-frequency carrier wave for wireless
transmission.
In 1870, John Tyndall demonstrated that short bursts of light pulses may also be
used in place of digital electronic pulses to transmit information. He further showed
that these light pulses may be sent through cables made of very fine glass fiber to
long distances without much loss. However, actual use of optical fiber communica-
tion started only in 1927. At present optical fiber communication is fast replacing
conventional metal core cable transmission for terrestrial communication. Optical
fiber communication at present is also being used for medical purposes, like in
endoscopes, for computer networking and in civil and military avionics.
https://doi.org/10.1007/978-3-031-32084-2_7
348 7 Optical Fiber Communication
7.2 Advantages of Optical Fiber Communication
For terrestrial transmission, optical fiber communication using light signals may be
compared with communication done using digitalised electronic signal via copper
core coaxial cables. Some of the important advantages of optical fiber communication
are listed below.
1. Increased information carrying capacity (large bandwidth): Optical fiber
provides high signal bandwidth resulting in significantly greater information
carrying capacity. In a way bandwidth may be compared with the diameter of
a pipe; if some fluid is transported through a pipe, pipe with bigger diameter
will carry more fluid. Moreover, if fluid is filled in a container, the container will
get filled faster with the pipe of larger diameter. Similarly, an optical fiber has
the capacity of holding and down (or up) loading of significantly larger amount
of data in comparison with a metal core coaxial cable. Optical fibers may be
classified as single mode (SM) and multimode (MM) types. Typical bandwidth
for (MM) type fiber is between 200 and 600 MHz-km, and it is more than 10 GHz
for (SM) fiber. The bandwidth for metal core cables ranges from 10 to 25 MHz-km
only.
2. Immunity from electromagnetic interference: In metal core cable transmis-
sion signal in the form of electric/electronic voltage pulse is used. Electronic
pulses in nearby cables interfere with each other (referred as cross-talking).
Electronic pulses also get affected by stray electromagnetic fields, particularly
high-frequency radio waves. No such interference of stray electromagnetic fields
is felt by the light pulses travelling through optical fibers. Therefore, transmission
is secure and faithful.
3. Low loss of pulse intensity: Designing of optical fiber is such that a light pulse
propagated through the fiber does not lose intensity or power, and this is achieved
by guiding the pulse through the fiber using total internal reflection. On the other
hand, electrical pulses inherently radiate energy and lose intensity as they travel
through metal core cable. No significant loss in light pulse intensity allows the
use of optical fiber cables of much higher lengths before a repeater station (an
intermediate station where pulses are re-strengthen) is installed. In comparison,
more repeater stations after shorter cable lengths are required for electrical signal
transmission using metal core cables. As a result of recent technological advances
in fabricating optical fiber, light can be guided through 1 km of fiber with an
intensity loss of as low as 0.16 dB (≈ 3.6%).
4. Reduction in size, weight and cost: Since optical fiber is made of silica glass,
optical fiber cables are light weight, small in diameter and cost substantially less
than a coaxial copper core cable. Normally, the diameter of an optical fiber is 1/
8 of the diameter of coaxial copper cable.
5. Enhanced security: Since leakage of light pulse from fiber cable is totally absent
and since light pulse propagated through fiber does not radiate, it is almost
impossible to detect the presence of underground fiber cable by scanning at
overground. Hence these fiber cables are very secure. In comparison, electronic
7.3 Basics of Optical Fiber Communication 349
pulse travelling through coaxial cable radiates electromagnetic energy that may
be detected by overground instruments, and their location may be detected rather
easily making them unsecure.
6. No grounding, no spark hazard: Since electric voltages are neither transmitted
nor develop during light pulse transmission through fiber cable, problems associ-
ated with sparking and with proper grounding do not arise in optical fiber trans-
mission. Optical fiber transmission is also immune to the potential difference
between the transmitting and receiving stations.
SAQ: In spite of so many advantages of optical fiber communication, there is one
big limitation. Can you point out the limitation?
7.3 Basics of Optical Fiber Communication
Optical fibers are mostly made from fused silica glass, as thin as human hair, and
designed to propagate light waves along their length using the principle of total
internal reflection. Light wave entering from one end of the fiber undergoes successive
total internal reflections from side walls and travels down the length of the fiber along
a zigzag path. At each total internal reflection almost all energy of the light wave is
reflected back into the fiber; however, only a negligibly small fraction of light pulse
intensity or power may escape from the side wall. As such the incident light wave
propagates through the fiber almost undiminished in power up to long fiber lengths.
Construction details of a typical optical fiber are shown in Fig. 7.1. A typical fiber
has a central core, generally of fused silica glass, of diameter ≈ 50 µm which is
surrounded by cladding. The refractive index ‘n’ of core material is always greater
than the refractive index n1 of cladding material. The overall diameter of cladding is
of the order of 125–200 µm. Total internal reflection of the light wave incident on
the fiber takes place at the boundary of core of higher refractive index and cladding
of lower refractive index. Silicon coating is provided over cladding, which improves
the quality of light transmission through the core. A tight-fitting buffer jacket,
generally made of plastic, surrounds Silicon coating. The main purpose of buffer
jacket is to protect the core and other coverings from absorbing moisture. Moisture
(–OH molecule) has very high absorption coefficient for light waves, and therefore,
attempt is made to totally eliminate the presence of moisture from the fiber assembly.
Buffer jacket is surrounded by a strength member, which provides mechanical
strength to the fiber cable. The whole assembly is covered by an outer jacket made
of polyurethane. This arrangement of different layers of shielding ensures that the
fiber cable is not damaged during pulling, stretching, bending, rolling, etc.
SAQ: What is the purpose of buffer jacket and strength member in a fiber cable?
Fig. 7.1 Construction

details of optical fiber
7.3.1 Optical Fiber Materials
Optical materials that are used in the fabrication of optical fiber must have following
properties:
(a) Material of the core and cladding must be transparent for the light wave lengths
used for transmitting optical signals in form of light pulses.
(b) Materials used for fabricating core and cladding should be such that their refrac-
tive indices may be manipulated by adding some impurities. As a matter of fact
the refractive index of core material should be slightly higher than that of the
cladding, which may be done by adding some impurities in core fiber material
or cladding material.
(c) Materials used for making both the core and the cladding must be such that they
may be drawn in long, very thin and flexible fibers.
Both glass (silica or silicate) and transparent plastic are good candidates for core
and cladding materials; glass has the edge of transmitting optical signals without
much attenuation to long fiber lengths as compared to plastic, while plastic has the
advantage of more mechanical strength as compared to glass. Therefore, depending
on the requirement, optical fiber cables of both glass and plastics are being used for
transmitting information in form of light signals.
In most glass fibers, silica with refractive index 1.458 (for light of wavelength
850 nm) is used as the basic material. Core material of higher refractive index may
be produced from this basic silica by adding small amounts of impurities like P2 O5 ,
GeO2 or B2 O3 , while silica itself may be used as cladding.
A new class of glasses, called ‘halide glasses’, has been recently developed which
has halide group (fluorine, chlorine, bromine, etc.) anions imbedded in them. These
glasses are found to transmit optical signals of infrared light with much smaller
loss over long distances. These glasses are used for core material, while cladding
material of lower refractive index is fabricated by replacing zirconium fluoride (ZrF4 )
by Hafnium fluoride (HaF4 ) in halide glasses.
Optical fiber cables with glass core and cladding made of plastic (Silicon resins)
have also been used for long-term applications. A popular cladding made of plastic
polymer is ethylene propylene with refractive index 1.338 (for IR wavelength 850 nm)
which is often used to cover the plastic core.
Plastic core and plastic cladding fiber cables are also in use, though they generally
have higher loss of signal strength over long transmission lengths. Methyl methacry-
lates (n = 1.49) cladding and polystyrenes core (n = 1.6) fiber cables are used for
their durability and rigidness.
7.3.2 Frequently Used Wavelengths in Optical Transmission
Even in the best optical fiber light is gradually attenuated as it travels through the fiber.
The linear attenuation coefficient or attenuation value is expressed in units of dB/km,
i.e. decibel per kilometre. Attenuation coefficient depends on the wavelength of light
used, and out of several causes of attenuation one important cause is the absorption
of optical signal light in water molecules (OH), that are invariably present in fiber
cables. Water molecules are imbedded in fused silica glass, being absorbed at the
time of manufacturing.
Linear attenuation as a function of wavelength of light is shown in Fig. 7.2. Three
natural dips in attenuation coefficient occur for wavelengths of 850 nm, 1300 nm
and 1550 nm; all the three wavelengths are in infrared region (IR) of light spectrum.
Light pulses of these three wavelengths are, therefore, often used for optical fiber
transmission. A big advantage of optical fiber is that three different data sets, each
coded in one these three wavelengths may be simultaneously transmitted by the
same fiber without any interference from each other. Simultaneous transmission of
multiple data sets by the same carrier is called multiplexing.
Other causes of signal attenuation will be discussed later in this section.
SAQ: What happens to the energy contained in a light photon when it is absorbed
by some molecule, like (OH) molecule?
7.3.3 Principle of Total Internal Reflection
Transmission of light wave through optical fiber without substantial loss of power has
been made possible by the principle of total internal reflection. Light waves may travel
through vacuum with characteristic high speed of 299,792,458 m/s ≈ 3 × 108 m/s,
denoted by ‘c’. However, in different mediums light travel with different speeds,
all smaller than ‘c’ the speed of light in vacuum. For example, the speed of light
in diamond is only 1.23881181 × 108 m/s and in pure water only 2.25407883 ×
108 m/s. The ratio (speed of light in vacuum/speed of light in medium) is called the
refractive index of the medium and is often denoted by either μ or n. Refractive index
of diamond is, therefore, 2.42 and of water 1.33. A medium that has a larger value
of refractive index as compared to another medium is called denser medium, and
Fig. 7.2 Linear attenuation coefficient as a function of light wavelength
light travels slower in a denser medium. The phenomenon of total internal reflection
occurs when light travels from a denser medium to a less dense (or rarer) medium.
Figure 7.3 shows four different situations of the passage of a light ray from denser
to a rarer medium. As shown in Fig. 7.2a, a ray of light in denser medium when
incident at the boundary of denser to rarer medium refracts in to the rarer medium.
The refracted ray shifts away from the normal; i.e. the angle of refraction in rarer
medium is larger than angle of incidence in denser medium. Further the angle of
refraction in rarer medium increases with the increase of the angle of incidence in
denser medium (Fig. 7.3b). On increasing the angle of incidence in denser medium
a situation is reached such that for angle of incidence of θC (in denser medium)
the refracted ray travels along the boundary separating the two media, as shown in
Fig. 7.3c. The angle of incidence θC for which the refracted ray travels along the
boundary is called the critical angle. On further increasing the angle of incidence
to values larger than the critical angle θC (Fig. 7.3d), the incident light ray does not
travel in the rarer medium, and instead it is reflected back in the denser medium. The
boundary separating the two media works as a mirror for incident angles larger than
critical angle θC (in denser medium). This phenomenon of reflection of light back
into the denser medium at the boundary of the two media (for incidence angle > θC )
is called total internal reflection. From Snell’s law it follows that
sin θc n1 n1 n
or θc = sin−1
1
= or sin θc =
sin 90 n n n
In ideal situation of total internal reflection no part of the energy contained in
the incident light ray is either transmitted to the rarer medium or is absorbed at the
Fig. 7.3 Passage of light ray from denser to rarer medium
boundary of the two media. Thus in total internal reflection a light ray is reflected
back in denser medium without any loss of its power or intensity.
Optical fiber is like a cylindrical waveguide made of low loss material, such that
the light incident on the core is guided through it (the core) by successive total
internal reflections from the core-cladding boundary. Only those incident rays that
hit core-cladding boundary with angles of incidence greater than critical angle θC
suffer total internal reflection and are guided further into the core without any loss
of intensity; other rays that hit the boundary with smaller angles of incidence may
refract into the cladding material and are lost.
7.3.4 Types of Fibers
Depending on the values of refractive indices of the core and the cladding, optical
fibers may be classified into three types.
(a) Single-mode step-index fiber
Figure 7.4 shows a single-mode step-index optical fiber. Let us understand what is
meant by the mode. In simple words, mode means the different paths of transmission
of the light signal through the core of the fiber. In a single-mode fiber, the diameter
of the core is very small so that all incident light rays after undergoing total internal
reflections from the core-cladding boundary travel (almost) along the axis of the core.
As such the transmission path for light rays through the core is only one, along the axis
of the cylindrical core. Hence such fibers are called single-mode fibers. Step index
means that the refractive index n of the core material is same over all sections of the
core, and the refractive index n1 of the cladding is smaller than n but remains constant
all over the cladding material. There is a step change in the value of the refractive
Fig. 7.4 Single-mode

step-index fiber
index at the core-cladding boundary, hence the name step-index fiber. Typical values
of the ratio of core diameter (2a) and the cladding diameter (2b) (both 2a and 2b in
micrometre µm) are: 8/125 and 50/125 for single-mode step-index fibers. If n and
n1 , respectively, denote the refractive indices of the core and the cladding, then the
fractional refractive index change Δ may be given as,
(n − n 1 )
Δ= (7.1)
n
Obviously, Δ is very small. Usually, the value of core refractive index n ranges
from 1.44 to 1.46, depending on the wavelength of the light, and the value of Δ lies
typically between 0.001 and 0.02.
(b) Multimode step-index fiber
A multimode step-index fiber is shown in Fig. 7.5. Typical diameter ratio (2a/2b) for
multimode fibers is 85/125 and 100/140. Since the diameter of the core is sufficiently
large, light pulse incident at the input end of the core may travel through several
different paths, along the axis, at many possible different angles such that after total
internal reflection at core-cladding boundary they all are guided through the core
and travel along different paths. Each of these paths represents a different mode
of transmission. In a step-index fiber, the refractive index n of the core and n1 of
cladding has fixed values. Refractive index of core is uniform throughout the core,
and similarly, the refractive index of cladding is uniform throughout the cladding;
there is a step change in the value of refractive index at the core-cladding boundary.
The problem with multimode step-index fiber is that light pulse or signal trans-
mitted through different modes (paths) travels different distances in the core. Since
the refractive index within the core has same value (n) in all parts of the core, light
signals travelling through different modes take different times to reach the terminal
point of the core. As a result, at the terminal point several images of the incident
light signal will be formed, each displaced with respect to the other by a very small
time interval. Thus the time definition of the incident signal will not remain sharp at
the terminal end. In an ideal situation images of the incident light signal travelling
through different modes along the fiber core should reach the terminal point at the
Fig. 7.5 Multimode step-index fiber
same instant, completely overlap each other to produce a sharp and intense image.
In a multimode step-index fiber final signal image is blurred.
(c) Multimode graded-index fiber
Figure 7.6 shows a multimode graded-index fiber. In multimode step-index fiber

the refractive index of core material was same in all parts of the core; as a result
incident light signal travelled different distances for different modes (paths) but with
same speed (as refractive index is same). This resulted in time dispersion in signals
reaching the end point from different modes. To overcome this problem, in graded-
index fiber the refractive index varies from the axis of the core towards its boundary,
which means that n is a function of n(r) the radius of the core; refractive index of the
core material has a maximum value at the core axis and decreases towards the core
boundary. However, the refractive index at the boundary of the core is still larger than
the refractive index of the cladding; hence total internal reflection at core-cladding
boundary does take place.
Fig. 7.6 Multimode graded-index fiber

Gradual decrease in the refractive index of core material from axis towards the core
periphery helps in reducing or eliminating time dispersion between signals travelling
through different modes. To understand the working of a graded-index multimode
fiber it may be recalled that light travels faster in medium of lower refractive index
as compared to the medium of higher refractive index. Incident light signals take
peripheral paths and have longer path lengths travel faster (because of the lower
value of refractive index), as compared to signals that travel via axial modes (shorter
paths lying near the core axis) through the medium of higher refractive index. In
nutshell, signals that travel through longer paths pass through medium of lower
index and therefore, travel faster, while signals that go through shorter paths travel in
medium of higher index and, therefore, move with lower speed. The net result is that
all signals, travelling through different modes in the core, reach the terminal point at
the same instant.
Light rays that are guided in the core suffer phase change at each total internal
reflection at core-cladding boundary. Some of the rays that are guided through the core
are in opposite phase and interfere destructively, annihilating each other, while those
rays that were in phase strengthen each other on undergoing constructive interference.
Hence, though several modes may be allowed but only some of them really take place.
7.3.5 Rays Guided Through Fiber
A light ray is guided by total internal reflection within the fiber core if its angle
of
incidence at the core-cladding boundary is greater than critical angle ∠θc = sin−1 nn1
and remains so as the ray undergoes total internal reflection from core-cladding
boundary again and again proceeding ahead.
7.3.6 Meridional and Skewed Rays
Rays in planes passing through the core (or fiber) axis are called meridional rays.
One such plane that passes through the core axis is shown in Fig. 7.7, and typical path
of a guided meridional ray is shown by ABCD. It may be noted that a meridional
plane always intersects the cylindrical core-cladding boundary at 90°. Therefore,
meridional rays intersect the fiber axis and are reflected in the same plane without any
change in their angle of incidence θ . These rays are guided if their angle i with the core

axis is smaller than the compliment of the critical angle θC [= π2 − θC = cos−1 nn1 ].
Since θC is usually small, meridional rays are approximately paraxial.
Any general ray is identified by its plane of incidence, a plane parallel to the
fiber axis and passing (or containing) the ray and the angle θ that the ray makes in
incidence plane. The incident ray after total internal reflection (called skewed ray) is
confined in a plane that is normal to the core-cladding boundary and is specified by
Fig. 7.7 Typical path of a guided meridional ray
Fig. 7.8 Path of a skewed ray through fiber core
the distance R by which the normal plane is offset from the core axis and the angle
ϕ that the normal plane makes with the incident plane, as shown in Fig. 7.8.
Guided paths in fiber core for meridional and skewed rays are shown in Fig. 7.9.
As may be seen in this figure, the meridional rays move in a plane, while the skewed
rays follow a helical path confined between two cylindrical shells.
SAQ: Explain why meridional rays are nearly parallel to the fiber axis?
7.3.7 Acceptance Angle
An optical fiber consists of a core of refractive index ‘n’ surrounded by a cladding

of slightly lower refractive index ‘n1 ’. Only those light rays that are incident on core
at such angles that they suffer total internal reflection on hitting the core-cladding
boundary are trapped within the core and are guided through it (the core). Further,
Fig. 7.9 Guided paths of

meridional and skewed rays
in a graded-index fiber
light rays impinging on the core-cladding boundary at an angle greater than the
critical angle undergo total internal reflection. In the following we will attempt to
find the value of the maximum angle of incidence for a ray of light that enters the core
from air (refractive index 1) and undergoes total internal reflection at core-cladding
interface. Acceptance angle is defined as the maximum angle of incidence at the
interface of air and core media for which the light ray enters the core and travels
along the boundary of the core and cladding.
Figure 7.10 shows a light ray incident on fiber core (refractive index n) from air
(refractive index 1). It may be noted in the figure that axis of cylindrical core acts
as normal to the air–core interface and that angle of incidence ∠i and the angle of
refraction ∠r are are related with each other by Snell’s law,
sin i n
= =n (7.2)
sin r 1
Condition for total internal reflection at point B in the figure is,
n1
sin ϕ ≥ (7.3)
n
But ∠ϕ = ∠(90 − r ), putting this value in Eq. (7.3), one gets;

n1 n1
sin ϕ ≥ or sin(90 − r ) ≥
n n
or
n1
cos r ≥
n
or
Fig. 7.10 Passage of a light ray incident on the fiber core from air
1/2 n 1
1 − sin2 r ≥
n
or
n 2
1
sin2 r ≤ 1 −
n
or
n 2 1/2
1
sin r ≤ 1 − (7.4)
n
But from Eq. (7.2), sin r = sin i

n
, substituting this Eq. (7.4), one gets;
or
n 2 1/2
sin i 1
≤ 1−
n n
or
1/2
sin i ≤ n 2 − n 21 (7.5)
Now for the case ∠i = ∠i acc , such that,

1/2
sin i acc = n 2 − n 21 (7.6)
The refracted ray will travel along the boundary of core-cladding interface. The
acceptance angle is, therefore, given by,
Fig. 7.11 Passage of light

rays incident on air–core
interface at different angles
of incidence
1/2
∠i acc = sin−1 n 2 − n 21 (7.7)
Rays incident from air to the fiber core making (with core axis) angles smaller
than acceptance angle ∠i acc will undergo total internal reflection at the core-
cladding boundary and will be guided in the core. A ray that is incident on the
air–core interface making an angle (with core axis) larger than the acceptance angle
will not undergo total internal reflection at the core-cladding boundary and will
refract into cladding. This is shown in Fig. 7.11, where ray AOBC incident at air–
core interface with incidence angle iacc after refraction at core-cladding boundary
travels along the boundary. Ray DOHK that makes angle of incidence larger than
acceptance angle (iacc ) is transmitted to the cladding after suffering refractopn at
core-cladding interface. Ray EOFG that is incident at air–core interface with angle
of incidence smaller than the acceptance angle undergoes total internal reflection at
core-cladding interface and is guided through thre fiber core.
7.3.8 Numerical Aperture (NA)

1/2
The quantity n 2 − n 21 is called numerical aperture (NA). It follows from Eq.
(7.6) that larger the value of NA, larger will be the angle of acceptance, and hence,
higher the light-gathering capacity. The numerical aperture, therefore, is a measure
of light-gathering capacity of the fiber. The magnitude of the numerical aperture may
also be written as,

1/2 (n − n 1 ) 1/2
NA = n 2 − n 21 = [(n + n 1 )(n − n 1 )]1/2 ∼
= 2n n.
n
or
(n − n 1 )
NA ∼
= n(2Δ)1/2 where Δ = (7.8)
n
Fig. 7.12 Acceptance cone and path of guided ray
As shown in Fig. 7.12 all rays incident at the air–core interface within the cone of
half angle i acc will fall at core-cladding boundary with angle of incidence larger than
the critical angle and will suffer total internal reflection and will be guided. On the
other hand, rays with angle of incidence at air–core interface larger than the angle of
acceptance i acc will suffer refraction at the core-cladding interface and will be lost
in cladding. Larger the opening of acceptance cone, more light rays may be guided,
which means that the light-gathering power of the fiber will be large.
SAQ: What is the physical significance of acceptance angle?
Solved Example SE7.1 Assuming that the core of a fiber is made of glass of refrac-
tive index 1.46 and is surrounded by a cladding of refractive index 1.40, calculate the
critical angle θc , numerical aperture NA, angle of acceptance and fractional refractive
index change Δ. What will happen to the critical angle, NA and angle of acceptance
if cladding is replaced by air of refractive index 1?
Solution In the first part of the question it is given that: n = 1.46 and n1 = 1.40.

Therefore ∠θc = sin−1 nn1 = sin−1 1.40 1.46
= sin−1 0.959 = 81.70◦ .
And θc = (90 − 81.70) = 8.30◦ .
1/2
Further, numerical apperture NA = n 2 − n 21 = 0.414.
And acceptence angle ∠i acc = sin−1 NA = sin−1 0.414 = 27.17◦
n − n1 1.46 − 1.40
Δ= = = 0.041
n 1.46
In second part of the problem it is asked to calculate the same quantities when n1
= 1.

−1 n 1 −1 1.00
∠θc = sin = sin = sin−1 0.685 = 48.03◦
n 1.46
And θc = (90 − 48.03) = 41.97◦ .

1/2 /
Also, numerical apperture NA = n 2 − n 21 = (1.46)2 − 1 = 1.0
And acceptence angle ∠i acc = sin−1 NA = sin−1 1 = 90◦ .
It may be observed that by replacing the cladding of index 1.40 by air of index 1.0,
i.e. by increasing the difference (n − n1 ) the angle of acceptance increases which in
turn will increase the light-gathering power of the fiber. It may, however, be shown
that with the increase of acceptance angle, the number of modes also increases
rapidly which is not good. It is for this reason that the difference between refractive
indices of core and cladding is kept small so that the number of modes does not
increase to a very large value.
7.3.9 The V Parameter
Light rays are electromagnetic waves with cross electric and magnetic fiels, and the
propogation of guided wave in the core medium o f a fiber may be studied following
the variation in electric and magnetic field vectors subject to Maxwell’s equations
and boundary conditions imposed by core-cladding indices. There are some special
solutions of these equations, called modes, each of which has a distinct propogation
constant, a characteristic field distribution in the transfer plane and two independent
polarisation states. In the analysis of guided electromagnetic waves a parameter
called fiber parameter or V parameter evolves which may be given as,
a a
V = 2π (NA) = 2π n(2Δ)1/2 (7.9)
λ λ
Here a is the radius of the core, λ the wavelenght of light, n the refractive index
of core material and Δ = n−nn
1
. V is an important fiber parameter that governs the
number of modes and their propogation constants.
The maximum number of modes N m supported by a step-index fiber is given by;
1 2
Nm = V (7.10)
2
However, for V < 2.405, the step-index fiber can support only one mode and is
called a single-mode step-index fiber (SMF). The wavelenght of light corresponding
to V = 2.405 is called the cut-off wavelength of the fiber and is often denoted by λc .
2πa
λc = n(2Δ)1/2 (7.11)
2.405
In case of graded-index fiber which may support large number of modes and has
large value of V, the maximum number of modes is given by;
1 2
Nm = V (7.12)
4
7.3.10 Attenuation and Dispersion of Optical Signal
There are two important parameters that are associated with any optical signal; they
are (i) power of the signal, often denoted by P, and (ii) spectral spread, i.e. the signal
has a specific wavelength λ with a spectral spread of Δλ. As the optical signal travel
sthrough the fiber its power diminishes, which is referred as signal attenuation. Since
different wavelength components of the optical signal travel with different speeds in
the fiber, the temporal spread of the signal becomes broadened, and this is referred as
dispersion. Attenuation and dispersion limit the performance of optical fiber medium
as a data transmission channel. Attenuation limits the magnitude of the optical power
transmitted, while dispersion limits the rate of data transmission.
(A) Attenuation
Experimentally it has been observed that the power of an optical signal transmitted
through optical fiber diminishes exponentially with distance on account of absorp-
tion in fiber medium and scattering. Attenuation is generally measured in terms of
the power transmission ratio and distance travelled by the signal. The attenuation
coefficient α in units of dB/km is given as;
1 Pinc
α= 10 log10 (7.13)
L PL
Here, L is the distance (generally in km) travelled by the optical signal in fiber, Pinc
the incident power of the signal at L = 0 and PL the power of the signal after trav-
elling a distance L in the fiber. Quantity PPincL is called the power transmission ratio.
The graphical relation between attenuation coefficient α and the power transmission
ratio is shown in Fig. 7.13. It may be observed in Fig. 7.13 that power transmission
ratio becomes 0.5 (signal power becomes half of its original value) when attenua-
tion coefficient has the value 3 dB. Similarly, attenuation of 10 dB corresponds to
power transmission ratio of 0.1 and 20 dB of transmission ratio of 0.01. It may be
observed that while attenuation coefficients are additive, power transmission ratios
are multiplicative.
Signal attenuation in optical fibers may be divided into two components: (i) due
to intrinsic causes that are due to the intrinsic properties of the fiber material and (ii)
extrinsic that are due to impurities, etc. not associated with the fiber material.
(i) Intrinsic causes of attenuation:
Two major causes of signal attenuation while travelling through the fiber are:
(a) Absorbtion in fiber core material and scattering. In optical fibers made of fused
silica glass core, absoption of optical signal strongly depends on the signal
wavelength. Fused silica glass has two strong absoption bands: one in ultraviolet
(UV) region (λ between 0.6 and 0.9 µm) and the other in infrared (IR) region
Fig. 7.13 Graph showing

dependence of attenuation
coefficient on power
transmission ratio
(λ between 1.6 and 1.9 µm). Between these two wavelength limits there is a
natural window for wavelengths in the region of nea- infrared region (λ between
1.0 and 1.6 µm) where there is no inherent absoption by fused silica glass.
Absorption of a given wavelength of light occurs when atoms and/or
molecules of the medium (SiO2 , fused silica) have either vibrational, rotational,
molecular or electronic excited states that exactly match with the energy of
the light wavelength. A light signal or photon of wavelength λ has an energy
ε = hc λ
and if the atom/molecule of the medium has some excited states that
differ exactly by ε amount of energy, the photon gets absorbed shifting the
system to the higher state of excitation. The near-infrared absorption band in
fused silica glass is due to vibrational bands, while the ultraviolet absorptrion
band arises because of the molecular and electronic excitations.
(b) Scattering of optical signal wavelengths by core medium is another major cause
of signal attenuation. The random localised variations of the molecular posi-
tions in silica glass of fiber core create random inhomogenities of the refrective
index and act as minute scattering centres. These inhomogenity centres scatter
optical signal light, in a way similar to the scattering of sun rays by dust parti-
cles in atmosphere. This type of scattering is called Rayleigh scattering and is
characterised by 1/λ4 law, which means that scattering is inversely proportional
to the fourth power of wavelength. Lower wavelengths are scattered more as
compared to the longer wavelengths. The wavelength window for which there
is no inherent attenuation of optical signal in fused silica glass is bounded by
Rayleigh scattering on the shorter wavelength side and by infrared absorption
on the long wavelength side.
(ii) Extrinsic causes of attenuation
Impurities, mostly of metallic ions and dissolved water vapours (OH ion) in silica
glass, cause attenuation of optical signals. Modern-day technology of making fused
silica glass has made it possible to remove almost 100% of all metallic ion impurities,
however. Water vapours dissolved in glass still cause absoption of optical signals. It
is for this reason that optical signals of some specific wavelengths, for which (OH)
ion absorptrion and Rayleigh scattering are minimum and intrinsic absorptions are
also small, are used. Bends and other physical distortions in fiber cables also cause
scattering/attenuation of signals. Further, attenuation of optical signal is least in step-
index single-mode fibers and is considerabley higher in multimode graded-index
cables.
(B) Dispersion
In optical fiber communication information coded in short-duration optical pulses is

transmitted from the transmitter end of the fiber cable network to the receiver end.
When a short-duration optical pulse travels through an optical fiber, the duration of
optical pulse may got stretched and its shape may get modified. Streching of pulse
duration and shape modifications depend on the distance travelled by the pulse in
the fiber. These changes in temporal and shape profiles of the incident optical pulse
while passing through optical fiber are referred as dispersion. For example, if the
time duration of the pulse at the transmitting end is say, t 1 , and at the other receiving
end of the optical cable it becomes t 2 (t 2 > t 1 ), then for reliable communication
next optical pulse from the transmitter end must be sent at least after a time t 2 from
the end of the previous optical pulse. Larger the value of t 2 , slower will be the
speed of communication. In other words larger the value of dispersion of the pulse,
slower becomes the speed of communication. Thus dispersion adversely effects the
speed of optical fibe communication. Figure 7.14 shows the dispersion produced by
a multimodal cable when a single light pulse is transmitted through it.
In order to understand the process of dispertion, it is required to know that an
actual optical pulse is a wave packet, composed of a spectrum of components of
different wavelengths each travelling in a medium with different group velocities.
This spread of group velocities of the wave packet is called the width of the wave
packet. The group velocity v of an optical pulse travelling in a dispersive medium of
refractive index n is given as,
c dn
v= , where N = n − λ (7.14)
N dλ
Fig. 7.14 Dispertion produced by a multimodal fiber cable

There are four main causes of signal dispersion in fibers. They are: (i) mate-
rial dispersion, (ii) modal dispersion, (iii) waveguide dispersion and (iv) nonlinear
dispersion.
(i) Matertial dispersion
Glass is a dispersive medium which essentially means that its refractive index is a
function of wavelength. Since an optical pulse is a wave packet, different components
waves of the wave packet travel with different group velocities in glass. It can be
shown that the temporal width of an optical pulse of spectral width σλ (in nm, 10−9 m)
after travelling a distance L becomes σr which is given by;
στ = |Dλ |σλ L (7.15)
where
λ d2 n
Dλ = − (7.16)
c dλ2
Dλ is called material dispersion coefficient. Generally, L is measure in km (kilo-
metre, 103 m), σλ in nm (nenometre, 10−9 m) and σr in ps (picco second = 10−12 s);
therefore, the units of dispersion coefficient are ps/km-nm. Wavelength depen-
dance of material dispersion coefficient for fused silica glass as a function of signal
wavelength is shown in Fig. 7.15.
As may be observed in this figure, coefficient has a negative value for wavelengths
below 1.3 µm. Wave packets of long wavelength travel faster than wave packets
of shorter
ps wavelength.
At λ = 1.3 µm D(1.3) = 0; forλ = 1.55 µm D(1.55) =
ps
+17 km−nm ; and for λ = 0.87 µm D(0.87) = −80 km−nm .
(ii) Modal dispersion
As the name suggests, modal dispersion takes place in multimodal transmission
cables. The reason for this dispersion is the fact that group velocities of signal through
Fig. 7.15 Material

dispersion coefficient for
fused silica glass as a
function of the wavelength
different modes are different. A single optical impuse of light (optical signal) entering
a M-mode fiber at z = 0 spreads into M signals, each travelling through the fiber
with a different group velocity. As a result, signals from different modes reach the
receiver end L (km) distance away after different time delays (see Fig. 7.14). If one
denotes by τk the time delay of the signal through the kth mode, then τk = vLk , where
vk is the group velocity of the signal through kth mode. If vmin and vmax represent
respectively the minimum and the maximum values of the group velocities, then
(L/vmin ) − (L/vmax ) will be the time spread of the optical pulse at receiver end.
The rms value of spread over all pulse widths at receiver end may be given by
σr = 21 [(L/vmin )−(L/vmax )]. σr is called the response time of the fiber. It can be
shown that in a step-index multimode fiber vmin = vmax (1 − Δ), hence,

1 L 1 1
στ = [(L/vmin )−(L/vmax )] = −
2 2 vmax (1 − Δ) vmax
L
= (1 − Δ)−1 − 1 (7.17)
2vmax
And since (1 − Δ)−1 ∼

= (1 + Δ), above expression reduces to,
L
στ = [Δ] (7.18)
2vmax
Response time or modal dispersion in graded-index fiber is much smaller as

group velocities for different modes are almost equalised. It may be shown that for
a greaded-index fiber the response time is given by,
2
L Δ LΔ2
στ = = (7.19)
2vmax 2 4vmax
(iii) Waveguide dispersion
Optical signal has electric and magnetic fields associated with it. Field distribution
in fiber depends on the ratio (a/λ), where a is the radius and λ the wavelength of the
light signal. Because of this dependance, group velocities of different modes differ
from each other, even when material dispersion is neglected. Dispersion produced as
a result of difference in modal group velocies on account of radius/wavelength ratio
is called waveguide dispersion. In analogy with Eq. (7.15), the waveguide response
may be written as,
στ = |Dw |σλ L (7.20)
where the coeffecient for waveguide dispersion Dw is given by,

1 d2 β
Dw = − V2 2 (7.21)
2π c dV
Here, V is the V parameter of the fiber and β the propogation constant of the incident
signal.
(iv) Nonlinear dispersion
Another type of dispersion occurs when the intensity of the input light signal is suffi-
ciently high. It occurs because for high-intensity light signal the refractive indices
become intensity dependent. High-intensity parts of the signal pulse undergo phase
shifts that are different from the low-intensity part. As a result the frequencies of
the signal get shifted by different amounts, which result in the modification of
group velocities and change in pulse shape. Under some special conditions, material
dispersion may counter balance nonlinear dispersion, so that the optical pulse travels
through the fiber without altering its temporal profile. The wave in such a condition
(when there is no change in its temporal profile) is called a soliton.
Figure 7.16a–c show dispersions produced by different types of fibers to a sharp
input optical pulse. It may be observed in these figures that the output pulse is shifted
in time with respect to the input pulse, and in most cases its shape has also changed
as a result of dispersion.
Fig. 7.16 Dispersion produced by different types of fibers

7.4 Components of Optical Fiber Network Link 369
SAQ: What is meant by dispertion of the signal in optical fiber transmission?

SAQ: Optical signal transmitted through a fiber reaches the destination after a time
delay which is proportional to the length of the fiber cable. Does this time
delay is a part of dispersion?
7.4 Components of Optical Fiber Network Link
In any information transmission the first step is to convert the given information
into digital data in form of electrical signals. For optical fiber communication the
informastion which is already available in digital electric/electronic form needs to
be converted into the form of optical (light) pulses. Essentially there are four main
components of a fiber network: (i) optical transmitter, (ii) optical fiber cable, (iii)
optical connectors and (iv) optical receiver. Figure 7.17 shows the simplistic layout
of a optical fiber network.
(i) Optical transmitter

It is a device that converts electrical signals into optical signals.Three types
of transmitters, depending on requirements, are in use. Light emitting diode
(LED) is the simplest transmitter that is often used in multimode applications.
Optical signal generated by a LED has largest spectral width that carries the
least amount of bandwidth.
Laser source which produces intense coherrent optical signal using the prin-
cipal of light amplification by simulated emission of radiations has smallest
spectral width and the largest bandwidth. Vertical cavity surface emitting laser
source (VCSEL) has spectral width less than that of LED.
Wavelength of the emitted optical signal depends on the material of the
transmitter device. As already mentioned, most of the time three wavelengths,
all in infrared region, 850 nm, 1300 nm and 1550 nm, are employed on account
of no absoption windows.
(ii) Optical connector
Fig. 7.17 Schematic layout of optical fiber network

Connectors are required to optically couple fiber to the the transmitter, to

the receiver and one fiber cable to another fiber cable to increase the length of
the cable. An optical connector terminates the optical fiber inside a ceramic
ferrule, using epoxy glue to hold the fiber in place. The advantage of connector
is that they may be connected (mated) or disconnected (unmated) at any time
and several number of times without damaging the connector. There are two
important types of connectors that are often used to link to some other device.
(a) Physical contact connector uses fiber in a tightly tolerance ceramic
ferrule, which allows easy handling and protects the fiber from phys-
ical damage. Fiber ends are highly polished, and when two connectors are
joined together, a sleeve over them ensures alignment of the two fibers
and a spring system keep ferrules tightly pressed so that ferrules remain
aligned and in constant contact. Physical contact connectors are regged,
reuseable, easy to clean and cost-effective. The loss of optical signal power
due to the use of physical contact connector (called insetion loss) is small,
of the order of 0.3 dB.
Figure 7.18 shows constructional details of a single fiber physical
contact connector. Multiple fiber physical contact connectors are also
available in the market.
(b) Expanded beam connectors use a lense at the exit end of each fiber
to widen and collimate the signal light. There is a air gap between two
connectors in this configuration as shown in Fig. 7.19.
Expanded beam connectors are less pron to dust and particle contam-
ination but are much susceptible to moisture and liquid contaminations.
Alignment of optical system is also difficult, and further the signal power
loss is also high of the order of 0.8–2.5 dB.
(iii) Fiber cable
Details of various types of optical fiber cables have already been dealt with in
previous sections.
(iv) Optical receiver
Fig. 7.18 Design details of single fiber physical contact connector

7.4 Components of Optical Fiber Network Link 371
Fig. 7.19 Expanded beam optical connector
Optical receiver receives transmitted information in the form of light pulses

and converts it in to corresponding electrical pulses. Most optical receivers use
photodiodes to convert optical pulses into electrical pulses. Two types of photo-
diodes p-type-intrinsic-n-type (PIN) diodes and avalanche photodiode (APD)
are often used in optical receivers. Special requirement of optical recevier is
that the photodiode must have high-frequency response since optical pulses
reaching diode are at high rate.
Photodiodes convert optical pulse into either a current pulse or a voltage
pulse. The current mode, where the output pulse is in the form of a current pulse,
is preferred as the current produed in photodiode output is linearly proportional
to the intensity of the incident optical pulse. A current pulse at the output can
always be converted into voltage pulse by passing it through a resistance. On
the other hand, in voltage output mode the voltage produced at the output
of a photodiode is logarithmically proportional to the intensity of the input
photopulse.
PIN photodiodes are slightly different in constraction from a normal pn junc-
tion photodiode. In a normal junction diode depletion region which behaves as
a capacitance is very narrow, and hence the capacity of depletion layer is large.
As a result of large capacitance, the frequency response of a normal junction
diode is slow. In a PIN diode the thickness of the depletion region is increased
by sandwitching a layer of intrinsic material between the p- and n-ends. Also
the p- and n-sides are heavyly doped. As a result of increased effective thickness
of depletion layer, the capacitance of the layer get reduced and the PIN diode
becomes capable to respond to high signal frequencies. Heavy doping of the
two sides further increases the strength of the electric field across the depletion
layer which in turn increases the speed of charge carriers through the region,
enhansing response to high frequencies. Figure 7.20 shows the layout of a PIN
photodiode. SiO2 coating is done over that area of the diode which has to be
masked from light, and Silicon nitride is painted to protect the interior from
oxydation, moisture and to reduce the reflection of light. The PIN photodiode
is kept reverse biased. Light signal reaches the depletion and intrinsic region
Fig. 7.20 Layout of a PIN photodiode
after crossing a thin p+ layer and produces electron hole pairs which are swept;
holes towards the p-side which is biased by negative potential and electrons
towards the n-side biased by positive potential. An electric current signal is
thus generated whenever a light signal shines the PIN diode.
Avalanche photodiode operates in reverse bias mode, with high reverse bias
voltage that may be up to 100 V. Optical (light) signal falling on depletion
region of a PAD produces charge carriers, holes and electrons which get accel-
erated under the high reverse bias and produce secondary charge carriers by
avalanche breakdown of depletion medium. As a result of secondary carrier
production high current flows through the photodiode generating a current pulse
corresponding to each light signal.
SAQ: What is the advantage of using a LASER-based transmitter in optical fiber

transmission?
SAQ: Among the physical contact and the expanded beam connectors which, in
your opinion, is better and why?
SAQ: What is the advantage of inserting a layer of intrinsic semiconductor between
the p- and n-type materials in a PIN diode?
SAQ: Is it advantageous to use photodiodes with large reverse bias voltage for
detecting light signal? Explain your answer.
7.5 Applications of Optical Fiber Transmission 373
7.5 Applications of Optical Fiber Transmission
Optical fiber transmission finds applications in different fields, on account of its

lightweight, high bandwidth capacity, high reliablity, no leakage, high safety and
immunity to electromagnetic and radio frequency (RF) interference.
(i) It is extensively used in avionics on both military and commercial aircraft
systems. Optical fiber communication is also used in those environments where
there are large number of heavy electrical motors in operation that generate
large electromagnetic back ground. Communication with metal core cables
in such environment easily picks up electromagnetic noise and distorts the
electrical signal used for communication and on the other hand communi-
cation done using optical fiber employs light pulses and fiber cables, and
the electromagnetic noise does not effect the optical signal and the signal
transmission.
(ii) Optical fiber communication is, therefore, used in drive control links for drilling
rigs and command and control of deep mining systems.
(iii) Fiber optics is also used in data communication and telecommunication due
to its ability to transmit high bandwidth over longer distances as compared to
electrical signal communication via copper core cables.
Solved Example SE7.2 Calculate the response time of a single-mode fiber that has
material dispersion coefficient Dλ of value −80 ps/km-nm for wavelength λ = 0.87
µm and the width of the signal (σ λ ) as 50 nm. Assume that no other dispersion except
material dispersion is present.
Solution Response time σ τ which is a measure of the broadning of the incident
optical pulse, for material dispersion, is given by;
στ = |Dλ |σλ L .
Substituting the values of Dλ and σλ given in the problem, one gets;

| |
| ps | 4000 ps 4 ns
στ = |−80 −nm|50 nm L = = L
km km km
From the expression above, one may infer that the rate of signal dispersion per
kilometre is 4 ns. It means that if an impulse light signal travels 100 km in the single-
mode fiber, then its width will become 4 ns × 100 = 4 × 10−9 × 100 s = 0.4 ×
10−6 s = 0.4 µs (Fig. 7.21).
Problem
P7.1 A photosignal of wavelength λ = 0.85 µm is guided through a fused silica

fiber with core refractive index 1.452, fractional refractive index Δ = 0.01
and core radius a = 25 µm. Determine the values of numerical aperture NA,
V parameter and approximate number of modes.
Fig. 7.21 Time response of an impulse optical pulse through a 100 km single-mode fiber
ANS: NA = 0.205, V = 37.9 and approximate number of modes = 585

P7.2 A step-index fiber has the radius a = 5 µm, core refractive index n = 1.45
and fractional refractive index change Δ = 0.002. Determine the shortest
wavelength of optical signal for which the fiber will allow single-mode
transmission.
ANS: 1.198 µm
P7.3 A step-index fiber is made with a core of index 1.52, diameter of 29 µm
and fractional index change Δ = 0.0007. It is operated with optical signal of
wavelength 1.3 µm. Find the value of V parameter for the fiober and possible
nmaximum number of modes it can sustain.
ANS: V = 4.049, maximum number of modes = 8
SA7.1 List the advantages and the drawback of optical fiber communication.
SA7.2 List important causes of signal attenuation in fiber communication.
SA7.3 What is dispesion of signal and how does it effect optical fiber communica-
tion?
SA7.4 Give a brief construction details of an optical fiber and explain the purpose
of each component.
SA7.5 What is the difference between a graded-index and a step-index multimode
fiber? Which is better and why?
SA7.6 What are the essential components of an optical fiber link? Give brief
description of each component.
SA7.7 Define acceptance angle, numerical aperture (NA) and V parameter of a fiber
and give physical significance of each.

Note: Some of the following questions may have more than one correct alternatives.
All correct alternatives must be marked for a complete answer.
MC7.1 Signal despersion coefficient Dλ has units
7.5 Applications of Optical Fiber Transmission 375
ps µs pm km−nm
(a) km−nm
(b) km−nm
(c) km−nm
(d) ps
ANS: (a)
MC7.2 Acceptance angle θacc is given by,
(a) sin−1 (NA) (b) sin−1 (2Δ) (c) sin−1 n(2Δ)1/2 (d) sin−1 n(2Δ)2
ANS: (a), (c)
MC7.3 Maridional rays are contained in a plane that
(a) Is at a fixed distance from fiber axis (b) passes through fiber axie (c) is
normal to fiber axis (d) makes an angle equal to the critical angle with fiber
axis
ANS: (b)
MC7.4 Bandwidth of a communication channel is a measure of the
(a) Time taken to transmit 100 kb data over 1 km (b) transmission capacity
of the channel (c) amount of data that may be up or down loaded per unit
time (d) change in temporal priofile of the transmitted signal
ANS: (b), (c)
MC7.5 In optical fiber transmission dispersion is a measure of the change in signal’s
(a) Power profile (b) frequency profile (c) shape profile (d) temporal profile
ANS: (d)
MC7.6 LASER source-based optical transmitter has
(a) Smallest spectral width (b) large bandwidth (c) least dispersion (d) least
attenuation
ANS: (a), (b)
MC7.7 The power of an optical signal becomes one-tenth of its original value
after travelling a distance of 100 km in an optical fiber. The attenuation
coefficient for the link in dB/km is
(a) 0.001 (b) 0.01 (c) 0.1 (d) 1.0
ANS: (c)
MC7.8 The V parameter for a 25 µm radius step-index fiber is 10. The maximum
number of modes supported by the fiber is
(a) 10 (b) 20 (c) 40 (d) 50
LA7.1 Describe in detail the components and their use in an optical fiber link.
Discuss the advantages of optical fiber communication and its applications.
LA7.2 Give details of the guided propogation of an optical pulse through a fiber
explaining the importance of acceptance angle, numerical apperture and V
parameter.
LA7.3 Discuss in detail the attenuation of optical signal while guided through a
fiber. What is dispersion and how does it adversly affect signal propogation?
LA7.4 Discuss different types of optical fibers and their merits. What is meant by
‘bandwidth’ and how is it related to the modes of fiber propogation?
Chapter 8
Laser Technology and Its Applications
Objective
Physics behind the working of LASER sources, their classification and some of their
important applications will be discussed in this chapter. After reading this chapter the
reader will be able to appreciate and understand special properties of LASER sources,
the processes of population inversion, induced emission and light amplification using
induced emission of radiations. He will also be able to follow why laser sources are
used in some special applications.
8.1 Introduction
LASER is acronym (short form) for ‘Light Amplification by Stimulated Emission of

Radiations’. Laser technology is used to fabricate special sources of light that delivers
intense, collimated, coherent and monochromatic light beams. Laser phenomenon
can be explained only by quantum mechanical treatment. Visible light, which human
eye can detect, is a small part of a much broader electromagnetic radiation (EM)
spectrum. In order to understand the working of a laser source it will be required
to know in some details about EM spectrum and interaction of EM radiation with
matter.
8.2 Electromagnetic Radiations
Electromagnetic (EM) radiations are emitted when a charged body is accelerated or

decelerated, when magnitudes of charges vary with time and also when the strength
of a magnetic field changes with time. Visible light is a very small part of a vast
electromagnetic radiation spectrum shown in Fig. 8.1.
https://doi.org/10.1007/978-3-031-32084-2_8
378 8 Laser Technology and Its Applications
Fig. 8.1 Spectrum of electromagnetic (EM) radiations
Broadly speaking, complete EM spectrum may be divided into two parts; visible
light and invisible light. Different parts of EM spectrum may be identified by their
specific names; like Gamma rays, X-rays, microwaves, radio waves, etc. As is shown
in the figure, different components of EM spectrum have different frequency or
wavelength ranges, but all EM waves travel with the same velocity c (= 3 × 108 m/
s) in vacuum. The frequency ν and the wavelength λ of an EM wave in vacuum
are related with the expression ν = c/λ. Since radiations transport energy from one
place to another, EM radiations carry with them energy. Quantum mechanics (QM),
the mechanics appropriate for microscopic systems, tells that energy carried by EM
radiations is in the form of small energy packets. Energy packets or energy quanta of
EM radiations are called Photon. The energy contents of a photon of an EM wave
of frequency ν are E = hν; here h is a constant called Planck’s constant with value
4.1357 × 10–15 eV s or 8.62 × 10–34 J s in SI units. As may be seen, Plank’s constant
is very small and, therefore, energy content of photon is very small. Order of energy
(in eV) contained in different components of EM waves is given at the top in Fig. 8.1.
Each component of EM radiations is of great use for humanity; X-rays are used to
take photographs of bones, etc., gamma rays for treating cancer patients, microwaves
for making ovens for cooking, ultraviolet light for sterilising surgical instruments and
killing bacteria, visible light for observing objects, infrared for heating and so on. In
order to use EM radiations, it is required to develop sources for different components
of EM radiations. Interaction of EM radiations with matter plays important role in
developing light sources.
8.3 Interaction of Electromagnetic Radiation with Matter 379
8.3 Interaction of Electromagnetic Radiation with Matter
Interaction of EM radiations with matter may be divided into two steps that compete
with each other (a) absorption and (b) spontaneous emission.
(a) Absorption
Materials are made up of atoms and molecules. When all atoms present in a material
are identical, the material is called an element. In the other case, when material
is made of different types of atoms, it may be called a salt, alloy, composite, etc.
depending on the nature of binding between atoms and molecules.
Each individual atom of a material, according the quantum theory, may exist in
one of the several quantised discrete energy states. Out of these energy states, state
with lowest energy is called the ground state and other states of higher energies,
the excited states. Energy states are not equidistant in energy from each other, the
energy separation between successive states decreases with the increase of energy.
In absence of any external radiations (at absolute zero temperature) all atoms of
the given material stay in their ground or lowest energy state which is most stable.
However, if the temperature of the material specimen is raised to some higher value T
K, some atoms absorb thermal energy from surroundings and shift to excited states.
At any given absolute temperature T Kelvin, when system is in thermal equilibrium,
number N E of atoms in an excited state of energy E is given by;
NE = N0 e(− kT )
E
(8.1)
Here, N 0 and k are respectively the number of atoms in ground state and Boltzmann
constant. Since the factor e(− kT ) is always less than one, N E < N (<E) . It means that
E
in normal conditions of thermal equilibrium, the population of atoms in any excited

state is less than the population of atoms in the state with lower excitation and the
population of the ground state. A schematic diagram of atom population in different
energy states at T = 0 K and T > 0 K for a system of atoms is shown in Fig. 8.2.
SAQ: What is meant by thermal equilibrium? How can one identify if a system is
in thermal equilibrium?
Fig. 8.2 Population of

energy states of an atom at T
= 0 K and T > 0 K
Fig. 8.3 Schematic diagram

showing absorption of a
photon by an atom
Let us now consider any two successive excited states of the atom with energies
E 1 and E 2 having atom populations N 1 and N 2 respectively. It may be noted that E 2
> E 1 but N 2 < N 1 . Now suppose a beam of EM radiations of frequency ν is made to
fall on the specimen. This beam of EM radiation will be a bunch of photons each of
energy E p = hν. If E p matches with the energy difference [E 2 − E 1 ] between two
energy states of the atom, i.e. if
hν = E 2 − E 1 (8.2)
then some of the incident photons may be absorbed by atoms in state E 1 and will
shift to the next higher state of energy E 2 . Absorption of one photon by one atom in
lower energy state E n will reduce the number of incident photons by one, decrease
the population of lower energy state by 1 and increase the population of atoms in the
higher energy state by 1. Schematic diagram showing absorption of one photon by
an atom in the lower energy state E 1 and shifting to higher energy state E 2 is shown
in Fig. 8.3. The process in which atoms in a lower energy state are made to absorb
photons externally incident on them is either called absorption or more specifically
induced absorption.
The number of absorption events N abs in time Δt is proportional to the number
of atoms N 1 in state E 1 , number N of photons per unit volume in EM beam and time
Δt. Number of photons per unit volume of the incident beam may be expressed as
energy density Q = Nhν of incident photon beam. Therefore,
Nabs ∝ N1 .Q.Δt
or
Nabs = B12 N1 .Q.Δt (8.3a)
Here, proportionality constant B12 is the probability of photon absorption transition.

The photon absorption rate R Abs may be given as;
Nabs
R Abs = = B12 N1 .Q (8.3b)
Δt
8.3 Interaction of Electromagnetic Radiation with Matter 381
Absorption of photons by atoms lifts them from a lower energy state to a state of
higher energy. The process is called photoexcitation. However, atoms of a material
may be excited by several other means; for example the atoms of the filament of a
bulb get excited when an electric current is passed through the filament. The atoms
of a solid body, like a metal ball, go to excited states when it is heated. Therefore,
photoabsorption is one of the many ways of atomic excitation.
(b) Spontaneous emission or de-excitation
Atom excited to a state of higher energy cannot live in the excited state for long. In
normal case the mean life time τ of atomic excited states is quite short of the order of
10−8 –10−9 s. Mean life is a statistical parameter. Number of atoms in an excited state
decreases exponentially with time. According to the exponential law of spontaneous
decay, if N (t=0) is the number of excited atoms at initial time t = 0, then number of
excited atoms N t at a later time t is given as;
Nt = N(t=0) e− τ
t
(8.4)
Figure 8.4 shows how, say 100 atoms in an excited state decreases exponentially
with time. If τ is the mean life for this decay, then nearly 63 (63%) of original
population of atoms in excited state will lose their extra energy and got de-excited.
After one mean life only about 37 (37%) of original population of excited atoms will
survive. About 98% of excited atoms will revert back to lower energy state after a
time of the order of 5τ. Reverting back from an excited state to the lower energy
state of an excited atom is called de-excitation. Mean life of an excited atomic state
is a parameter that cannot be controlled by any physical or chemical process; it
is a property of the excited state which cannot be altered. Natural de-excitation is a
spontaneous process. An important observation that follows from the spontaneous de-
excitation curve is that atoms in a given excited state do not de-excite simultaneously
at the same instant; 98% of excited atoms de-excite in a time span of about five mean
lives.
Excited atoms may shake-off their extra energy and de-excite in several ways,
including the spontaneous emission of a photons. Emission of photons from the
excited atoms, without any external stimulation, is termed as spontaneous photon
de-excitation. In a way it is just the opposite of absorption. In normal course both
Fig. 8.4 Spontaneous

de-excitation of excited
atoms with time
photon absorption and spontaneous photon emission occur simultaneously when a

specimen is hit by an external beam of photons of appropriate energy.
In spontaneous photon emission, photons of same energy (hν) are emitted but
generally at different instants of time and in different directions. Figure 8.5a shows
two energy states, lower energy state with larger atom population and higher one with
less population of excited atoms. At a given instant one of the excited atoms in higher
energy state spontaneously de-excites, decreasing the population of higher energy
state by 1, increasing the population of lower energy state by 1 and emitting a photon
of energy E p = hν, in a random direction (see Fig. 8.5b). At some other instant another
excited atom spontaneously emits a photon of energy hν in some different random
direction as shown in frame (c) of Fig. 8.5. About 98% of excited atoms will de-excite
by photon emission in a period of time of the order of 5 times the mean life of the
excited state, and most of the emitted photons will be emitted at random moments,
moving in random directions. Except of their energies, there will be no co-relation
between spontaneously emitted photons. A light source based on spontaneous photon
emission may deliver photons of same energy but not a collimated and co-related
(coherent) beam of photons. Most of daily use light sources, like an electric bulb,
use spontaneous photon emission from excited atoms and deliver, incoherent and
uncollimated beam of light.
Coming back to spontaneous photon emission, the number of spontaneous tran-
sitions N sp taking place in time interval Δt will be proportional to the number of
excited atoms in the state of higher energy N 2 , and the time interval Δt, i.e.;
Fig. 8.5 a Populations of the lower and the excited states, b, c show the spontaneous photon
emission, d shows that de-excitation photons are generally emitted at different times and in different
directions
8.4 Einstein Prediction of Stimulated Emission 383
Nsp = A21 N2 .Δt (8.5)
And the spontaneous decay rate (N sp /Δt) is given by

sp
Spontaneous decay rate RD = A21 N2 (8.5a)
Here A21 is the probability of spontaneous photon de-excitation. It may be noted

that spontaneous decay rate does not depend on the energy density, Q that is on the
number of incident photons N.
SAQ: Suppose a collection of atoms at some temperature T is bombarded with
photons of appropriate energy so that some ground state atoms go to the higher
excited state by induced absorption, at the same time some excited atoms
might be undergoing spontaneous de-excitation. What will be the relation
between the induced absorption and spontaneous decay rates in the state of
thermal equilibrium?
(c) Radiationless de-excitation

An atomic or molecular excited state de-excites to a state of lower energy mostly
by spontaneous emission of photons; however, some times the excited state reverts
to a state of lower excitation without emitting any photons. Spontaneous decay of
an excited state without emitting photons is called radiationless de-excitation. In
case of radiationless decay, the excess energy is lost through what are called ‘relax-
ation processes’, that includes molecular dissociation, inelastic collisions, rotation,
vibration and other similar interactions. Exponential law of spontaneous decay holds
good for radiationless transitions also.
SAQ: According to Eq. (8.4), decay rate of excited atoms depends on the mean life
τ; however, the mean life τ does not explicitly appear in Eq. (8.5). How can
this be explained.
8.4 Einstein Prediction of Stimulated Emission
The origins of stimulate emission, the back bone of laser physics, may be traced
back to an idea from Albert Einstein in the formative years of quantum theory. In
1917 Einstein published a paper entitled ‘Zur Quantentheorie des Strahlung’ which
in English translates as ‘On the Quantum theory of Radiation’. In this paper Einstein
predicted stimulated emission of radiations in order to satisfactorily reproduce the
quantum mechanical expression of radiation energy density Q. Earlier, way back in
1900 an expression Q was derived by Planck.
Einstein argued that in case of an atomic system having two energy states E 2 and
E 1 (E 2 > E 1 ) with number of atoms N 2 and N 1 (N 2 < N 1 ) respectively, if bombarded
by a beam of N photons per unit volume, each of energy hν = (E 2 − E 1 ), the system
will be in thermal equilibrium only when the decay rate (by spontaneous emission
of photons) of state E 2 to state E 1 and the photon absorption rate of state E 1 to state
E 2 are equal. These transition rates are given respectively by Eqs. (8.5a) and (8.3b).
On equating the two rates one gets,
A21 N2 = B12 N1 .Q
or
A21 N2
Q= (8.5b)
B12 N1
E2 E1
However, from QM N2 = N0 e− kT and N1 = N0 e− kT .
− ( 2kT 1 )
E −E
Therefore, N2
N1
= e substituting this value of N 2 /N 1 in Eq. (8.5b) one
gets,

A21 − ( E2 −E1 )
Q= e kT
(8.5c)
B12
Earlier, another German scientist Max Planck carried out detailed study of black-
body radiations and in 1900 gave the following expression for the energy density of
radiations Q,
8π hν 3 1
Q= 3
hν (8.5d)
c e kT − 1
Einstein found that the two expressions for energy density of radiations, given
by Eqs. (8.5c) and (8.5d), do not agree. He therefore suggested that there may be
another factor that has been missed while driving Eq. (8.5c).
Einstein proposed that apart from spontaneous decay there may be a stimulated
or induced decay of atoms in excited state E 2 that might also contribute to the decay
process. He further assumed that the decay rate RDsti due to stimulated emission is
proportional to the number of atoms N 2 in state E 2 and also the energy density Q.
Thus,
RDsti ∝ N2 Q Or RDsti = B21 N2 Q (8.5e)
Here B21 is the constant of proportionality but different from A21 . Another important
point to note is that Einstein assumed that stimulated decay rate is not only propor-
tional to N 2 , the number of atoms in state E 2 but also to the energy density Q (=
N hν).
With the addition of another mode of decay the total decay rate of level E 2 is now
the sum of decay rates due to spontaneous emission and decay rate due to stimulated
emission. In thermal equilibrium the sum of the spontaneous and stimulated decay
rates must be equal to the absorption rate;
8.4 Einstein Prediction of Stimulated Emission 385
sp
RDsti + RD = R Abs
or
A21 N2
B21 N2 Q + A21 N2 = B12 Q N1 or Q =
(B12 N1 − B21 N2 )
or
A21
Q=
B12 NN21 − B21
( E2 −E1 )
Substituting the value of N1
N2
= 1
( E2 −E1 ) = e kT in above expression for Q,
−
e kT
one gets;
⎛ ⎞
A21 A21 ⎝ 1 ⎠ = A21 1
Q= ( E2 −E1 )
= ( E2 −E1 ) B12 (hν)
B12 e kT − B21 B21 B12
e kT − 1 B21
B21
e kT −1
B21
(8.5f)
8π hν 3 1
If one now compares Eq. (8.5f) with the expression of Q = c3 hν given
e kT −1
by Planck Eq. (8.5d), one gets
A21 8π hν 3 B12
= 3
and = 1 (or B12 = B21 ) (8.5g)
B21 c B21
Constants B12 , B21 and A21 are called Einstein’s coefficients and relations specified
by Eq. (8.5g) Einstein’s relations.
The ratio R of rates for spontaneous emission to stimulation emission may be
given as,
A21 N2 hν hν
R= = e kT − 1 ≈ e kT (8.5h)
B21 N2 Q
Equation (8.5h) tells that the ratio of spontaneous to simulated decay increases
exponentially with energy of photon. For example if one calculates ratio R for the
light of frequency 4.7 × 1014 emitted by ruby laser at room temperature (300 K) using
the values of h = 8.63 × 10−34 J s and Boltzmann constant k = 1.28 × 10−23 J/K; one
gets R = 3.7 × 1032 , it means that when there is thermal equilibrium spontaneous
decay will be 1032 times more probable than stimulated decay. Obviously, stimulated
decay in thermal equilibrium will be completely masked by spontaneous decay. In
order to enhance stimulated decay one has to go to a situation where the system is
not in thermal equilibrium. Further details to achieve such situation will be discussed
in next sections.
SAQ: In what respect the spontaneous de-excitation differs from stimulated de-
excitation?
Willis Lamb and R. C. Retherford were the first in 1947 to experimentally observe
stimulated emission. In 1950 Alfred Kastler gave the idea creating a non-equilibrium
state by population inversion for which he was awarded Nobel Prize of Physics in
1968.
8.5 Stimulated (or Induced) Emission of Photons
Let us consider two states of an atom with excitation energies E 1 and E 2 (E 2 > E 1 ).
Further, let there be N 2 excited atoms in state E 2 and almost no atom in state E 1 . It is
also assumed that conservation laws allow photon de-excitation of atoms in state E 2
to E 1 . In normal situation atoms from state E 2 will follow spontaneous exponential
decay and will de-excite to state E 1 emitting photons of energy (E 2 − E 1 ) in random
directions with random phases and states of polarisation.
However, a special phenomenon, called induced de-excitation or induced/
stimulated photon emission may occur if a photon of energy hν = (E 2 − E 1 ) is
made to hit excited atoms of state E 2 . The incident photon will immediately make
one excited atom of state E 2 to de-excite by emitting a photon of energy hν = (E 2 −
E 1 ). This photo-de-excitation of excited atom by another photon of same energy is
called induced or stimulated photon emission. In stimulated photon emission, the
incident photon simply induces an excited atom to emit a photon identical to the inci-
dent photon, without undergoing any change in itself. Induced de-excitation does not
depend on the mean life of the excited state E2 , and it occurs immediately as the
incident photon interacts with an excited atom in state E 2 . Further, the phase, state
of polarisation and direction of motion of the incident photon and the photon
produced by stimulation emission are same.
With stimulated photon emission by the incident photon, there are now two
photons of same energy, same phase, same state of polarisation and moving in
the same direction. The two photons each of energy (E 2 − E 1 ) now induce two
more excited atoms to de-excite by emitting two more photons. Thus a single inci-
dent photon produces a chain reaction in which the numbers of induced photons
multiply rapidly. As a result almost all excited atoms in state E 2 undergo stimu-
lated de-excitation at the same instant by emitting large number of photons that are
all moving in the same direction, have same energies, have same phases, and same
states of polarisation. It is important to note that stimulated emission of photons does
not depend on the mean life of the excited state, and it occurs almost instantaneously.
Frame (a) in Fig. 8.6 shows two energy states of a system with energies E 2
and E 1 (E 2 > E 1 ) with one excited atom in state E 2 and an empty state E 1, before
8.5 Stimulated (or Induced) Emission of Photons 387
some instant t. At instant of time t, a photon of energy E p = hν = (E 2 − E 1 ) is

incident on the system. The incident photon stimulates the excited atom which de-
excites by emitting an induced photon of same energy E p , same phase, same state
of polarisation and moving in the same direction at instant of time t. In frame (b) of
the figure, before instant t, there were many excited atoms in state E 2 , all of which
de-excite simultaneously at the same instant t when a photon of energy E p was
incident. A bunch of photons, all having same energies, phases, states of polarisation
and direction of motion is produced at the instant t by stimulated emission from de-
exciting atoms. The bunch of coherent photons, all moving in same direction in the
same part of space, may undergo constructive interference resulting in a very intense
beam of monoenergetic photons. Triggering of the process in which large number
of photons of same energy, direction of motion, phase and state of polarisation are
produced by a single incident (light) photon is called ‘light amplification by simulated
emission’ or the laser action. A light source that uses laser technology to produce
intense photon beam is called a leaser source of light (EM radiations). It may be
observed that stimulated emission initiated by a single incident photon results in the
production of large number of coherent photons. Thus it may be said that stimulated
photon emission results in producing photon gain. The medium that contains atoms/
molecules that may be excited and may be made to undergo stimulated emission is
called active medium, or gain medium. Active medium may be solid, liquid or gas.
Fig. 8.6 Stimulated emission of photons a when there is only one excited atom, b when there is
large number of atoms in excited state
The number of stimulated transitions N sti occurring in time interval Δt is given

as,
Nsti = B21 N2 QΔt (8.6)
Here, B21 is the probability of stimulated transition, N2 the number of excited atoms
in state of energy E 2 and Q the energy density of incident photon beam.
8.5.1 Population Inversion
While discussing the laser action above, it was assumed that the population of excited
atoms in state of higher energy E 2 is large while the lower energy state E 1 is empty.
This assumption that there were large number of excited atoms in state E 2 and either
no or very few atoms in lower energy state E 1 is physically not justified. Quantum
statistics tells that in thermal equilibrium, when a system is in steady state, the number
of excited atoms in a state of higher energy is always less than their number in a state
of lower energy as dictated by Eq. (8.1). The situation when number of excited atoms
is more in a state of higher energy than their number in a state of lower energy is
technically referred as population inversion. It is obvious that population inversion
is normally not possible, and if a system has population inversion then it is not in
thermal equilibrium. However, as will be seen, in some special cases, population
inversion can be achieved using some technical tricks.
Let us now discuss why population inversion is a prerequisite for viable laser
action. Let us assume that there are more atoms in the state E 1 of lower energy
and only few excited atoms in state E 2 of higher energy, i.e. there is no population
inversion. In this situation, if some stimulated photons are released by laser action
from the de-exciting atoms in state E 2 , they (stimulated photons) may be absorbed
immediately by atoms in state E 1 . Since the number of atoms in state E 1 is much larger
than that in state E 2 , the probability of photon absorption will be large; as a result
either negligible or no coherent photons will be available for making a light source,
which means that there will be no gain of photons in the system even after stimulated
emission. It is, therefore, necessary to have population inversion for having large
photon gain factor so as to have a viable laser technology based source.
Some important characteristics of simulated emission (laser action) are;
• It requires population inversion between two chosen energy states of the system.
• It may be initiated by a photon of energy hν equal to the energy difference between
the two states.
• Photons emitted in stimulated emission are identical in all respect to the incident
photon, they are coherent, having same energy (or frequency), same phase, same
state of polarisation and same direction of motion.
• The process of stimulated emission may be controlled from outside.
• In process of stimulation emission multiplication of photons takes place.
• Coherent photons may undergo constructive interference resulting in a highly

collimated, intense and monoenergetic beam of photons.
SAQ: Does the state of population inversion an equilibrium state?
8.5.2 Essential Requirements for Laser Action
There are some essential requirements for practical application of laser technique,
they are;
I. There should be two states of an atomic system, one of higher energy E2
and the other with lower energy E1 such that photon transition from state
E2 to state E1 is allowed. Quantum states of a system have several other good
quantum numbers, like spin angular momentum, parity, magnetic moment, etc.,
apart from its energy. Transition between two states of a system is governed
by some conservation laws. It is possible that photon transition between some
energy states is prohibited on account of some conservation law. However, for
lasing action to take place, photon transition from higher lasing state E 2 to the
lower lasing state E 1 must be allowed.
II. The state with higher energy E2 must have a mean life of at least 10−4 s
or more. As has already been discussed, population inversion is essential for
practical application of laser action. It means that the population of atoms in
state E 2 must be increased to a high value before triggering stimulated emission.
Normally, the highest population of atoms is in ground state, and it is required
to lift atoms from the ground state to the state of energy E 2 either directly or
via some other intermediate step. The process/technique of lifting in energy
atoms from a state of lower energy to the state of higher energy for simulated
emission is called pumping. Pumping takes time to create population inversion
to the desired level. If state E 2 is very short lived, with life time of the order of
10−8 s, atoms pumped to state E 2 will decay out by spontaneous emission before
sufficient degree of population inversion is achieved. Hence, it is required that
the mean life of state E 2 should be of the order of 10−4 s or larger. Nature has
provided atomic systems with excited states that have longer live times. Those
excited states that have live times larger than 10−8 s are called metastable states.
These metastable states are generally chosen as the upper state for lasing action.
III. Pumping is the process of supplying energy to atoms in a state of lower
energy so as to lift them to the state of higher energy to create and maintain
population inversion. Pumping essentially is a method of putting additional
energy in the active medium to produce and sustain population inversion.
Energy, to atoms in lower energy state (mostly the ground state), may be given
in many different forms, like via absorption of photons by atoms, via some
chemical reaction or through inelastic collisions of atoms, etc. Some details of
these pumping options will be discussed in the following.
8.5.3 Pumping
(a) Optical pumping
French Physicist Alfred Kastler in 1950 proposed a method to alter the relative popu-
lation of excited levels by optical irradiation of atoms in ground state. He visualised
a two-step process where in atoms in some level A are hit with a beam of photons of
an appropriate energy, which they absorb (called induced absorption) to move to the
excited state B. Some of the excited atoms in state B may revert back to level A and
some others to another excited state C (may be a metastable state). If de-excitation
rate of level C to level A is lower than the feeding rate from state B, population of
state C will increase with time at the cost of the population of A, see Fig. 8.7.
Optical pumping is often used in solid state lasers.
(b) Electric discharge or excitation by electrons
In case of optical pumping energy to atoms in lower energy state (generally the ground
state) is supplied by shining a beam of light (photons) on them. In case of pumping
by electric discharge, the required energy to atoms in ground state is supplied by
establishing a large potential difference across the system, which are mostly gases.
Under high electric field gas atoms/molecules get ionised emitting electrons. These
electrons get accelerated under the high electric field and collide with other atoms/
molecules to excite them to higher energies. Some of the excited atoms/molecules
may feed an intermediate state (exactly as in case of optical pumping) the population
of which may increase with time at the cost of the ground state population. This
method of pumping is used in gas lasers like argon laser.
(c) Inelastic atom–atom collision
This method of achieving population inversion is used when active medium is a

mixture of two gases, say A and B. An electric discharge is passed through the active
medium, which excites atoms of the gas with lower value of excitation potential
to their excited state, denoted by A*. Excited atoms A* collide with ground state
atoms B of the other gas that has higher ionisation potential. The inelastic collisions
Fig. 8.7 Producing

population inversion using
optical pumping
between excited atom A* with ground state atoms B, transfer excitation and a part
of kinetic energy to atoms B and excites it to a higher excitation state B*. Atoms in
excited state B* may partly de-excite to the ground state B or to some intermediate
state C. If the decay rate to the intermediate state C is larger than the decay rate to
the ground state, the population of intermediate state C will increase with time and
may become larger than the population of ground state B. He–Neon (He–Ne) laser
uses this mechanism of pumping.
(d) Thermal pumping
Source of energy that shifts ground state atoms/molecules to higher excited state in
case of optical pumping is photons of characteristic energy, and in cases of electric
discharge and inelastic atom–atom collisions the electrostatic field. Desired atomic/
molecular excited states in some cases may also be achieved by heating the active
laser medium. Heat energy may work as the energy source for pumping. Except that
heat energy becomes the source of excitation, the population inversion in this case
also is achieved the same way as it is in optical pumping.
(e) Chemical pumping
Some chemical reactions leave the product molecule or atom in an excited state,
while the ground state of the system (atom/molecule) is unstable or dissociative.
The energy released in chemical reaction is converted into the excitation energy.
Such chemical reactions automatically produce population inversion and are used
for making laser sources. Following chemical reactions are often used in laser sources
that may deliver power up to hundreds of watts.
Reaction Laser
H + Br2 → HBr∗ + Br HBr
F + H2 → HF∗ + H HF
F + D2 → DF∗ + D DF
CS + O → CO∗ + S CO
(f) Pumping based on direct conversion of electrical energy into light
This technique is used in fabricating solid state semiconductor diode lasers. These
diode lasers are compact with an active medium in the solid phase, which directly
converts electrical energy into laser radiations. The laser energy output of diode lasers
may be gainfully employed to further pump other solid state lasers.
8.5.4 Three and Four Level Lasing Schemes
Atomic and molecular systems have several excited states with at least one or two
of these as metastable states. In cases where the first excited state of the system
is a metastable state, lasing action may be achieved between the ground and the
metastable state, under suitable conditions.
Layout of a three-level lasing scheme is shown in Fig. 8.8. Atoms/molecules from
the ground state are pumped using the appropriate pumping method to the energy
level E 2 . Level E 2 undergoes spontaneous de-excitation both to the metastable state
E 1 with decay constant λ2 and to the ground state with decay constant λ1 . In case the
decay constant (number of decays per unit time) λ2 is much larger than λ1 , (λ2 >> λ1 )
and the mean life of the metastable state E 1 is sufficiently large, the population of state
E 1 will increase with pumping time at the cost of the population of the ground state.
Eventually, after some time of continuous pumping, population inversion between
the ground and the first excited state E 1 will be achieved. A photon of energy (E 1 −
E g ) may now trigger lasing action between the two states. In this scheme the first
excited state E 1 is the upper lasing level while the ground state the lower lasing level.
In three-level lasing scheme (Fig. 8.8) the ground state is populated directly from
energy level E 2 through spontaneous decay (decay constant λ1 ) and also from level
E 1 via lasing. As such continuous pumping at sufficiently rapid rate is essential for
maintaining population inversion. This drawback may be overcome in four-level
lasing as shown in Fig. 8.9.
In four-level lasing, atoms/molecules from the ground state are pumped to an
excited state E 3 which decays partially to the ground state E g with decay constant
λ1 and mostly to another excited state E 2 with large decay constant λ2 . State E 2
is metastable. There is another energy state E 1 which is fed by state E 2 , however,
as E 2 is metastable state spontaneous decay of E 2 to state E 1 is very weak. State
E 1 de-excites to ground state with decay constant λ3 . With appropriate pumping,
population of metastable state may be increased with respect to the population of
state E 1 and population inversion may be easily achieved and maintained since state
E 1 spontaneously decays to ground state. The metastable state E 2 and the state E 1
are respectively the upper and the lower lasing levels.
Fig. 8.8 Layout of a

three-level lasing scheme
Fig. 8.9 Layout of four-level lasing scheme
8.5.5 Optical Resonator or Laser Cavity
Once population inversion is established in a system, spontaneous de-excitation of

the higher lasing level to the lower lasing level produces the triggering photons
for lasing action. But the spontaneous de-excitation photons may be emitted in any
direction; therefore, laser light produced by the triggering of different spontaneous
decay photons will not be in one direction, instead will be oriented randomly. In
order to obtain an intense laser beam in only one direction it is required to select and
intensify laser radiations in only one direction. This is achieved by the use of optical
resonator or cavity resonator arrangement.
Laser cavity or optical resonator is generally an arrangement of few fully reflecting
and one partially reflecting mirrors surrounding the active medium. Multiple reflec-
tions of laser light between mirrors intensify and deliver an intense laser beam in a
specific direction.
An optical resonator is an arrangement of optical components, which allows
a beam of light to circulate in a closed path so that it traces its own path multiple
times, in order to increase the effective length of the media with the aim of large light
amplification. The process is similar to the process of positive feedback in electronic
amplifiers. As already mentioned, photons emitted in spontaneous de-excitation of
the upper lasing level works in initial stages as the seed photons for lasing action,
but later photons produced by lasing action and constrained to move in closed path
by optical components of the resonance cavity initiate stimulated emission and thus
amplify or multiply the number of laser photons in the closed path. Larger the path
length of laser photons in active media more will be the amplification factor. In this
way the active medium works as an amplifier and carry a certain gain factor.
(i) Gain coefficient of the active medium
Let us assume that there is a system with two energy states E 2 and E 1 (E 2 > E 1 )
with N 1 and N 2 number of atoms (N 1 > N 2 ) and photon transition between the two
states is allowed. Let the system be kept in some medium. Suppose that a beam of
monochromatic photons of frequency ν = (E2 −E h
1)
and intensity I0 is projected in
the medium. In normal situation when there is no population inversion, photons from
the incident beam will be absorbed by the atoms in the lower energy state as the beam
travels through the medium. As a result the intensity of the beam will exponentially
decrease with distance x travelled in the medium, i.e. Ix = I0 e−αx , the coefficient
of absorption α will be proportional to the difference (N 1 − N 2 ). Now suppose that
optical pumping is done and population inversion occurs resulting in laser action and
emission of laser photons of frequency ν. Under the changed circumstances N 2 is
now greater than N 1 and since α is proportion to (N 1 − N 2 ), it will become −α and
the intensity of the incident photon beam will be given by; Ix = I0 e+kx , where k (=
−α) is a constant that will be proportional to (N 2 − N 1 ). It means that the intensity
of the incident beam will increase as it will travel the active media. Here k is called
the gain coefficient of the medium. It can be shown that the gain coefficient k of the
medium is given as,
μ(N2 − N1 )hν B(2→1)

k= (8.6a)
c
where μ is the refractive index of the medium.
(ii) Threshold gain coefficient for lasing
In resonator cavity laser beam undergoes multiple oscillations in the cavity media
between the pair of mirrors to obtain large gain before leaving the cavity through
the partially polished mirror. Laser oscillations can only sustain in the active media
of the cavity if it attains at least unit gain after a round trip from mirror to mirror
and overcome the various losses in the cavity. Laser beam propagating through the
cavity medium undergoes scattering, etc. with medium atoms/molecules and lose
its intensity. If γ denotes the coefficient of beam intensity loss in the medium, the
overall gain coefficient of the beam will become (k − γ ).
Let us consider a laser beam of initial intensity I 0 that starts from mirror M 1
towards mirror M 2 as shown in Fig. 8.10. As the beam hits mirror M 2 its intensity
becomes I 1 given by I1 = I0 e(k−γ )L . The beam suffers reflection at mirror M 2 and if it
is assumed that the refractivity of the two mirrors is respectively R1 and R2 , intensity
of the beam after reflection at M 2 will become I2 = R2 I1 = R2 I0 e(k−γ )L . While
travelling from M 2 to M 1 through the active medium the beam intensity will change
and will have the magnitude I3 given by I3 = I2 e(k−γ )L = R2 I0 e(k−γ )L e(k−γ )L =
R2 I0 e2(k−γ )L at the instant when the beam hits mirror M 1 . On suffering another
reflection at M 1 the beam intensity will change to I4 = R1 I3 = R1 R2 I0 e2(k−γ )L .
Therefore, the net gain G in the intensity of laser beam in one round trip between
two mirrors of the cavity is given by G = II04 = R1 R2 e2(k−γ )L . The gain G of the
laser beam must either be 1 or larger than 1 for sustained oscillations of the beam in
the cavity. Therefore, the threshold condition for sustained oscillations of laser beam
in the cavity is given by
Fig. 8.10 Change in

intensity of laser beam in one
round trip between mirrors
G = R1 R2 e2(k−γ )L = 1
That gives e2(k−γ )L = R11R2 or (k − γ ) = 2L1

log R11R2 .
The threshold gain of the cavity medium (k − γ ) must satisfy the above condition
for sustained laser beam oscillations in the cavity. If laser beam gain per trip between
the two mirrors is less than the threshold gain, the beam oscillations will die out.
The gain factor for an active medium originates from the fact that excited atoms/
molecules produced by pumping are randomly distributed in the cavity medium at
different locations. Each excited atom is a source of laser photon when stimulated.
Larger the number density of excited atoms in the medium more will be the gain of
the medium.
SAQ: In what form the gain is distributed in active media?
(iii) Axial or longitudinal modes
The layout of a plane parallel resonator also called Fabry–Perot optical cavity, that
uses one fully reflecting and one partially reflecting plain mirrors around the active
medium is shown in Fig. 8.13. Two mirrors are held normal to the optical axis of
the medium. Spontaneous de-excitation photons trigger laser bunches that move in
random directions as shown in Fig. 8.13a. Laser bunches moving at some angle with
the optic axis are allowed to be lost but bunches moving along the optical axis stays
in the medium and undergo multiple reflections from the two mirrors. Photons of
bunches moving along the optical axis trigger stimulated emission from excited atoms
in the higher lasing level, thus generating additional laser bunches all moving along
the axis (see Fig. 8.13b). As a result of multiple reflections, the medium is filled with
large number of laser bunches all in phase but moving in opposite directions along
the optical axis. This happens only if the gain per trip of photon bunches between the
two mirrors is either equal or larger than the threshold gain. Laser bunches moving
in opposite directions produce stationary or standing waves. For standing wave, it
is required that the optical path length travelled by a wave between consecutive
reflections should be an integer multiple of wave length. If L is the distance between

the two mirrors, then
2L = mλm (8.7a)
where integer m may have values = 1, 2, 3 . . .

Equation (8.7a) tells that several laser bunches having slightly different wave-
lengths λ1 = 2L; λ2 = L , λ3 = 23 L . . . etc. may have sustained oscillations in the
cavity, subject to the condition of threshold gain. Expression (8.7a) in terms of photon
frequency may be written as,
c
νm = m (8.7b)
2L
here c is the velocity of light in vacuum
In case the active medium has refractive index μ, then the above expression gets
modified to
c/μ
νm = m (8.8)
2L
In principle laser light is monochromatic which means that all photons in a
bunch of laser light have exactly same energy (or frequency/wavelength) given by
(E 2 − E 1 ). However, this is not true, the fact is that a laser beam contains large
number of photons of exactly same energy but there are also photons, small in
number, that have energies slightly more and slightly less than the energy of the
most abundant photons. The intensity profile of a laser beam is like the one shown
in Fig. 8.11; as shown in the figure maximum number of photons are of frequency
ν max, but there are photons in smaller number, of frequencies different from ν max.
The spread of frequencies is called the laser line width and is denoted by Δw. Laser
line widths are generally of the order of 0.5 nm. Reasons for the width of laser beam
are discussed later in the section.
Some frequencies within the line width Δw satisfy condition laid down by Eq.
(8.8) and undergo sustained oscillations in resonance cavity. Other frequencies in
Fig. 8.11 Intensity profile

of a laser beam
the line width that do not satisfy sustain oscillation condition, die out. Final inten-
sity pattern of the laser beam emerging from the cavity through partially polished
mirror M 2 is shown in Fig. 8.12a. These frequencies; ν1 , ν2 , ν3 etc constitute axial
or longitudinal modes.
The difference Δν in the frequencies of any two successive modes is given as c/μ2L
from Eq. (8.8) as shown in Fig. 8.12b.
Fig. 8.12 a, b Axial modes and separation between modes
Fig. 8.13 a Fabry–Perot resonator showing bunches of laser photons moving in different directions.
Photon bunches are produced by lasing action initiated with spontaneous de-excitation photons
moving in random directions. b Laser bunches of photons moving along the optic axis of the active
material undergo multiple reflections at two mirrors and stimulate further lasing action adding
more coherent photons. Coherent photon waves strengthen each other by constructive interference
producing an intense laser beam that emerges from the partially reflecting mirror. c Laser bunches
moving along the optical axis but in opposite directions produce stationary waves of multiple
longitudinal modes
Fig. 8.14 Typical shapes of beam spot corresponding to different TEM configurations
It follows from Eq. (8.7a) that the number ‘m’ of possible axial modes for a given
laser beam may be calculated as, m = 2L λ
, where λ is the wavelength of maximum
emission. For example, in case of Ruby laser λ = 694.3 nm = 694.3 × 10−9 m and
2L
hence the number of possible axial modes in this case may be as large as 694.3×10 −9 . If
−2
the length L of cavity is say, 25 cm, then number of axial modes may be 2×25×10
694.3×10−9
=
72×10 . Further, Δν = 2L the frequency difference between two successive modes
4 c/μ
3×108
for L = 25 cm and μ = 1, comes out to be 2×25×10 −2 = 8.0 × 10 Hz.
8
The line width Δw for ruby laser is typically around 0.3 GHz (= 0.3 × 109 Hz or
0.53 nm). It may be observed here that for the case of ruby laser with cavity length
of 25 cm, Δν = 6 × 108 Hz is more than Δw = 0.3 × 109 Hz . This means that
only one axial mode will be sustained by the cavity (Fig. 8.13).
(iv) Transverse modes
Longitudinal or axial modes are generated by photons of frequencies different from

the frequency of the most abundant photons and moving along the axis of the mirror
system. Transverse modes are produced by photons moving at small angles with
respect to the mirror axis. Laser beam photons, only slightly off the mirror axis might
also produce standing waves and may undergo sustained oscillations in optical cavity,
giving rise to transverse modes. Transverse modes affect the shape of the laser beam
cross section and may be identified by looking at the beam cross section, when the
laser beam hits a surface. These modes, represented by TEM (short form for trans-
verse electromagnetic modes) are designated by putting appropriate subscript to the
short form TEM. Figure 8.14 shows shapes of laser beam spots and the corresponding
TEM symbols.
Four different types of optical cavities often used for achieving large signal
gain are shown below.
8.6 Special Characteristics of Laser Light 399
8.6 Special Characteristics of Laser Light
Laser light has three special properties that light emitted by an ordinary incan-
descent electric bulb or any such source do not have. These properties are: (i)
monochromaticity, (ii) coherence and (iii) directionality.
(i) Monochromaticity
Laser light photons are emitted when atoms/molecules from some excited state of
energy E 2 are stimulated to decay to a state of lower energy E 1 (see Fig. 8.9).
Since, the upper and lower lasing states are quantum mechanical states, they have
discrete and definite energies, in the present case E 2 and E 1 respectively. Therefore,
all the emitted laser photons have the same energy E p = (E 2 − E 1 ). This property
that all emitted laser photons have same energy (or wavelength or frequency) is
called monochromaticity. In contrast, light photons emitted by other sources like
incandescent lamps have a spectrum of wavelengths (energies), from infrared to
ultraviolet.
A detailed study, however, shows that though the energy or wavelength of all laser
photons is very nearly same, but there is some small spread in the wavelength (or
frequency) of laser photons also. There are three main reasons for this small spread
in wavelength of laser photons.
(a) Natural line width
Though it is said that quantum states E 2 and E 1 have discrete energies, but these
energies have some inherent uncertainties. These uncertainties in energy of excited
state arise from another fundamental law of quantum mechanics called the uncer-
tainty principle. According to the uncertainty principle, the uncertainty ΔE in energy
E of a quantum mechanical state and its mean life τ are related by the relation;
ΔE.τ ≈ h/2 (8.9)
Here, h is Planck’s constant. It follows from Eq. (8.9) that ΔE will be zero only when
mean life τ is infinite. Since the mean lives of atomic and molecular excited states
are not infinite, there is always some uncertainties ΔE 2 and ΔE 1 associated with
the energies of the upper and lower lasing levels. These uncertainties are referred as
natural line widths of energy levels. Uncertainties in energies of the upper and the
lower lasing levels produce spread in the energy E p of laser photons;
E p = (E 2 ∓ ΔE 2 ) − (E 1 ∓ ΔE 1 ) = (E 2 − E 1 ) ∓ ΔE pNat (8.10)
Spread ΔE pNat in laser photon energies is often called natural spread or natural
width of energy level.
(b) Doppler broadening
It is a common experience that the frequency of the whistle sound emitted by a
moving train appears changed when the train approaches the observer and when it
moves away from the observer. This change in frequency is called Doppler effect. A
similar change in frequency (or wavelength or energy) of laser photons occurs when
they are emitted by moving atoms/molecules in upper lasing state E 2 . According to
the kinetic theory of matter, any object at a temperature T > 0 K, has some kinetic
energy, as such atoms/molecules in upper lasing level E 2 are in state of random
motion and the frequency of the laser photon they will emit will be different if the
emitting atom/molecule is moving in the direction of the emitted photon or opposite
Dop
to it. Therefore, another uncertainty ΔE p in the energy of laser photons, called
Doppler uncertainty, is introduced by the random thermal motion of emitting atom/
molecule. Of course Doppler uncertainty will depend on the temperature of the laser
source.
(c) Recoil broadening
A photon of frequency ν carries a linear momentum P = hν c
. Let us assume that
an atom/molecule in upper level E 2 is stationary and emits a photon of frequency ν.
Since emitted photon has carried a linear momentum P, the emitting atom/molecule
must recoil back to conserve the linear momentum. The recoil of the atom will give
atom some velocity, say v and energy E atom = 1/2M atom v 2 . The question is who will
supply this energy E atom to the recoiling atom? The answer is the emitted photon. As
a matter of fact the photon will not be emitted with frequency ν, it will be emitted
with some lower frequency ν ' so that the energy difference h(ν − ν ' ) may take care
of atomic recoil. Actually, at some temperature T > 0 K, atoms/molecules in state
E 2 are not stationary but in random motion, different atoms/molecules will recoil
by different amounts and this will introduce additional spread ΔE pRec in the emitted
laser photons.
(d) Energy bands in solid state lasers

In case of laser sources where the active medium is solid, the energy level schemes
of atoms/molecules or ions are made up of energy bands rather than energy levels. It
happens because of the interaction between nearby atoms, which is quite appreciable
in solids as compared with liquids or gases, on account of close packing of atoms
in solids. Widths of both, the upper and the lower lasing bands introduce additional
uncertainty E pBand in the wavelength of emitted laser light, only in solid state lasers.
The net effect of all the above factors is to produce an overall uncertainty or
spread in the energy (frequency or wavelength) of emitted laser photons denoted by
Dop
E ptotal = E pNat ∓ E p ∓ E pRec ∓ E pBand . In spite of the above-mentioned causes the
overall uncertainty or spread in the wavelength of laser light Δλ is generally < 1 nm
(10 Å), while the spread of wavelength in ordinary light from incandescent lamp
may be as large as 10–1000 nm.
Uncertainty or the spread of any spectral line may be given either in terms of the
wavelength Δλ or in terms of the frequency Δν, and the two values are related by
the relation;
c c
Since λ = , hence dλ = − 2 dν
ν ν
Replacing dλ by Δλ and Δν by Δν one gets
c
Δλ = − Δν
ν2
Negative sign in the above expression simply tells that when λ increases ν
decreases.
(e) Laser groups from a system
In some special cases, two or three groups of monoenergetic laser light photons, each
with its own energy uncertainty, may be emitted from the same system.
Let us consider an atomic/ion/molecular system that has energy level structure
as shown in Fig. 8.15a. Atoms/ions of the system may be pumped from the ground
state E g by some suitable pumping mechanism to the excited state of energy E 5 ,
which spontaneously decays to the metastable state E 4 . Since E 4 is a metastable state
of comparatively long mean life, spontaneous decay of state E 5 accumulates large
number of atoms/molecules in energy state E 4 producing population inversion with

respect to the other three energy states E 3 , E 2 and E 1 , which undergo spontaneous
de-excitation to the ground state. Metastable state E 4 may spontaneously decay to the
three intermediate energy states E 3 , E 2 and E 1 , each with a certain branching ratio.
Lasing action may, therefore, take place between energy states E 4 and E 3 , between
E 4 and E 2 and between E 4 and E 1 emitting laser groups of wave lengths, say
hc hc hc
λ1 = ; λ2 = ; λ3 = (8.11)
(E 4 − E 3 ) (E 4 − E 2 ) (E 4 − E 1 )
The intensity of each laser group will depend on its branching ratio, and each
laser light will have its own wavelength uncertainty, ∓Δλ1 , ∓Δλ2 and ∓ Δλ3 ,
respectively. It is possible to select any one of the three laser lights by manipulating
or adding suitable devices in optical resonator. This way of selecting a particular
laser light of a given wavelength out of several groups (three in the present example)
is technically called tuning to the particular frequency or wavelength.
Figure 8.15b shows another energy level scheme of some atomic/ionic or molec-
ular system, which has got two metastable states E 4 and E 3 which may be fed by
the spontaneous decay of the higher energy state E 5 . The two metastable states
may decay, respectively, to states E 2 and E 1 . Appropriate pumping mechanism may
continuously populate state E 5 which in turn feeds metastable states to produce the
state of population inversion between E 4 and E 2 and E 3 and E 1 . Photons produced by
the spontaneous de-excitation of metastable states may initiate laser action resulting
in two groups of lasers with wave lengths λ1 ∓ Δλ1 and λ2 ∓ Δλ2 .
(ii) Coherence
In lasing action almost all excited atoms/molecules in the upper lasing level E 2
undergo stimulated de-excitation at the same instant releasing large number of
photons. Each of these photons is an electromagnetic wave; the special property
of these electromagnetic waves is that they are all in phase and are moving in same
direction. In-phase means that the troughs and crests of all these electromagnetic
Fig. 8.15 Emission of laser light of different wavelengths by a system

waves occur at the same instant, i.e.; the crests and troughs of different electromag-
netic waves fall on each other. The large number of in-phase electromagnetic waves
moving in the active medium in same direction interfere constructively giving rise
to a single light wave of very large amplitude. Since the intensity of a light wave is
proportional to the square of the amplitude, the constructive interference of photon
waves results in a very high intensity light beam.
(iii) Directionality
A laser source emits light only in one direction, in contrast to other light sources,
like electric bulb, light-emitting diode or tube light which emit light in all directions.
Resonance or optical cavity associated with a laser source is responsible for the
property of directionality of laser light. For example, in case of the optical cavity
that has two plain mirrors, one fully and one partially reflecting, held parallel to each
other; the direction of the laser light is defined by the direction of the normal to the
two mirrors, from full to partially reflecting mirror.
Generally the intensity profile of a laser beam has Gaussian distribution; intensity
is maximum, say I 0 at the centre of the beam and it falls off rapidly towards the edges
as shown in Fig. 8.16. The points where the intensity drops to 1/e of its central value
are called outer edge points. Distance 2w between the outer edge points is referred as
the minimum beam spot size. Directionality of laser beam is often defined in terms
of the full angle divergence ϕ, given as;
1.27λ
ϕ= (8.12)
2w
Here λ is the wavelength of the laser light
In a practical laser source light comes out from an aperture of diameter ‘d’ (which
is different from 2w) and spreads out because of diffraction. Assuming laser light to
be a plane wavefront, it may be shown that the light propagates as a parallel beam
2
for a distance of dλ , called the Rayleigh’s range, and there after it spreads because
of diffraction. The angular spread of the laser beam due to diffraction Δθ is given
by,
λ
Δθ = (≈ 10−5 or 10−6 rad) (8.13)
d
Fig. 8.16 Intensity profile

of a laser beam
Fig. 8.17 Laser light

diverging from a point source
In case of laser the angular spread Δθ is generally less than 0.01 mrad; meaning
there by that a laser beam spreads less than 0.01 mm for every 1 m of its propagation
distance. In contrast, the angular spread for light from an ordinary source may be
1 m in travel path of 1 m or more.
When light diverges from a point source (see Fig. 8.17) and r 1 and r 2 are the radii
of the beam at distances d 1 and d 2 , respectively, the angle of beam divergence in
radians is given as,
(r2 − r1 )
ϕ= (8.13a)
(d2 − d1 )
If ϕ is the angle of beam divergence of a laser beam then the area of beam spread
Aspr at a distance D from the source is given by,
Aspr = (Dϕ)2 (8.13b)
(iv) Irradiance
The irradiance (power per unit area) of a typical laser source is far greater than other
sources of light. It is largely due to the directionality and compactness of the laser
beam. One may take the example of an ordinary bulb of 10 W power, which spreads
its light uniformly in all directions, the irradiance I bulb at a distance of 1 m from the
bulb will Ibulb = 4π.1 2 m2 = 0.8 W/m . On the other hand, a typical laser light of
10 W 2
power output as small as 1 mW is concentrated in a radius of 1 mm at a distance of

1 m. The radiance of the laser light I laser is, therefore, Ilaser = π.(0.001)2 2 = 318 m2 .
0.001 W
m
W
(v) Focusability
An ordinary source of light must have an appreciable transverse extent so as to
produce a significant amount of light energy. The images of these non-laser sources
formed by lenses and mirrors have finite size which may be determined from the
laws of geometrical optics. The amount of light (energy) at the image position is
determined by the amount of light from the source intercepted by the lens or mirror.
As a result, in case of ordinary light sources large amount of light energy cannot be
focused at a small spot. In case of laser source, on the other hand, the transverse extent
of a laser source is very small which allows the lens or mirror to intercept almost all of
the power of the laser beam and focus it at a very small spot at image point. Moreover,
8.7 Classification of Laser Sources 405
because of the high directionality of laser beam, laser beam behaves as a bundle of
parallel light rays coming from a point source at infinity. As a consequence almost
total laser beam power gets focused at a very small spot at image position. Basically,
the size of the image spot is mostly governed by diffraction and lens aberrations and
can be as small as the wavelength of laser light. Property of focusability makes it
possible to deposit large amount of laser beam energy in a very small area to drill fine
holes in thick and hard sheets of matter. Because of the special property of focusing
in spots of the size of light wavelength, laser beam is used in ocular surgery where
target irradiance of ≈ 109 to 1012 W/cm2 is required.
Above-mentioned special characteristics, monochromaticity, coherence and
directionality, of laser light makes it extremely useful but also more hazardous.
One should never look directly on a laser beam because the highly collimated beam
may focus to nearly a microscopic dot on the retina of the eye, causing almost instant
damage to the retina. Coherence of laser light is important for observing interfer-
ence effects that has important applications in precision measurement of distances.
The branch of optics, interferometry, uses superimposition of coherent light to make
extremely fine measurements of very small distances (of the order of the wavelength
of light), surface irregularities or changes of refractive index. High intensity of laser
light, an outcome of coherence, deposits large amount of (heat) energy in a very
small area hit by the beam. The output power of a laser source may vary from a
few thousandths of a watt (in case of laser pointers) to several thousand watts. Laser
beams may be used as a cutting tool for thick metallic sheets or for welding of
metallic components. Laser light is also finding applications in medical field where
laser sources are used in equipment used for endoscopy and other surgical operations
including treating cancer cells.
SAQ: Which factor(s) determine the intensity of image formed by lens or mirror?
Why images formed using laser sources are very bright?
8.7 Classification of Laser Sources
Laser sources may be classified in several different ways;

(i) According to the physical state of the active medium:
(a) Solis state lasers (b) Gas lasers (c) Liquid lasers and (d) Semiconductor
lasers
(ii) According to the mode of operation:
(a) Pulsed and (b) Continuous wave (CW) lasers
(iii) According to other properties
(a) Gain of active medium (b) Efficiency (c) Power delivered and (d)
Applications.
In the following we will discuss laser sources according to the classification based
on the physical state of the gain medium or active medium.
8.7.1 Solid State Lasers
As the name suggests, laser sources in which the active medium is a solid material
are included in this class. Solid state lasers may be further divided into two types: (i)
doped insulator rod type and (ii) semiconductor diode laser.
(i) Doped insulator rod type lasers
Most of solid state laser sources use a rod of some crystalline insulating material
called the host which is doped with some ions, called the active ion. The first laser
ever made was a solid state laser which used a synthetic ruby (aluminium oxide
crystal) as host and chromium as an active species. Later, other solid state lasers
using sapphire as host and titanium atom as active species, etc. were also developed.
Generally, solid state lasers are named by prefixing the chemical symbol of the
active specie followed by the name of the host crystal, for example, Cr.ruby laser or
Ti.sapphire laser, Nd.YAG (Neodymium. Yttrium Aluminium Garnet, Y3 Al5 O12 )
laser, etc.
Commonly used host materials are, Y3 Al5 O12 (yttrium aluminium garnet, YAG),
Al2 O3 (sapphire), YVO4 (Yttrium orthovanadate),Y3 Sc2 Al3 O12 , Al2 BeO4 (alexan-
drite), Mg2 SiO4 (forsterite), etc. Silicate and phosphate glasses are also used as host
materials.
Light-emitting atoms, called active species, are embedded in the host rod. Active
ions may be of (a) transition metals like, Fe2 + , CO2 + , V2 + , Cr3 + , Ti3 + , Ni2 + , etc.;
(b) of rare earths like, Nd3 + , Pr3 + , Sm2 + , Pm3 + , Eu2,3 + , Dy3 + , Tb3 + , Ho3 + , Yb3 + ,
etc. and actinides such as U3 + .
The layout of most of the solid state laser sources is similar. Host crystal with
doped active ions is taken in the form of a rod, the dimensions of which may vary
from few cm to few tens of cm in length and a fraction of cm to few cm in diameter.
In most cases optical pumping is used to achieve population inversion. Since atomic/
molecular excited states in solids have band structure, optical pumping with a light
source that has broad spectrum covering the absorption band of laser material is most
efficient way of pumping. Both pulsed and continuous light sources may be used for
optical pumping, however, in case of continuous light source unabsorbed light energy
increases the temperature of the laser tube. Therefore, in such cases cooling of laser
cell is required. On the other hand pulsed light sources that produce flashes of light
of few micro- to millisecond duration with a gap of about the same duration are often
used to reduce heating problem. The active ion implanted host rod and the source of
pumping light/flash lamp(s) are generally kept in an elliptical glass tube, the inside
of which is polished for concentration of the light emitted by the flash lamp on to
the laser rod by reflection from polished surfaces. Several different arrangements of
mirrors of different curvatures, including parabolic, etc. are used to focus flash light
on the host material rod for maximum pumping speed (see Fig. 8.18). Laser sources
that use flash light are inherently pulsed laser sources as pumping and population
inversion occurs only for the duration of the flash light. The choice of flash lamp
depends on the energy of the active ion state(s) that are required to be populated
by pumping and, therefore, on the level scheme of the active ion. Halogen lamps
and high-pressure mercury discharge lamps are generally used as continuous light
sources for pumping, while low-pressure quartz or glass-sealed xenon or krypton
lamps are used as pulsed light sources.
Figure 8.19a shows a typical circuit that may be used to produce flashes of light.
As shown in the figure, energy storing capacitor is charged through a resistance R
to the maximum value V (few kV) of power supply voltage. A coil of few turns
called triggering coil, wrapped round the flash tube may be energised periodically
by a RF power supply. Triggering coil when energised produces ionisation of some
gas molecules in flash tube, thus reducing the electric resistance of the tube. Once
electric resistance of flash tube is reduced, the high voltage across the energy storing
capacitor produces an electric discharge in the flash tube emitting light. The duration
of the discharge is governed by the time constant of the LCR1 circuit, where R1 is
the equivalent resistance of the flash tube. The off-time of the flash lamp is governed
by the charging time of the capacitor to the power supply voltage through resistance
Fig. 8.18 Different geometrical arrangements of flash lamps and host laser rod in solid state lasers
Fig. 8.19 Electronic circuit for operating a flash lamp
R. Sometimes the flash tube is made in the form of a helix, and the lasing-doped
insulator rod is placed inside the helix for efficient pumping (see Fig. 8.19b). The
purpose of surrounding the light source and the lasing rod with mirror system and
focusing the light on rod is, maximise the flux of emitted light on the rod so that large
number of ions from lower energy state(s) may be pumped to higher energy state(s)
for population inversion.
Overall efficiency is an important parameter of a laser system. Overall efficiency is
defined as the ratio of the total power required to pump the laser to the optical output
power of the system. Typical efficiencies range from fractions of percent to 25%
or more. Many important laser sources have efficiencies as small as 0.05%, in such
cases the difference of between the input power and the output power is converted
into heat. Most laser sources, therefore, have a cooling system associated with them.
Laser systems with solid gain medium are typically cooled by surrounding the gain
medium in a cooling jacket; water or oil flows through the jacket to remove heat.
Low power sources may be cooled by forced air cooling, and very low power systems
may not need any cooling at all. Further details of Cr.ruby and Nd.YAG laser sources
are given in the following.
(a) Cr.ruby laser source
First ever successful laser source was made by Theodore Maiman on 16 May 1960
at Hughes Research Laboratories in California using a chromium doped synthetic
ruby rod. Present-day ruby laser has undergone several improvements since its initial
stage.
A rod of synthetic ruby crystal is made by doping a small amount (0.05% by
weight) of chromium oxide (Cr2 O3 ) in aluminium oxide (Al2 O3 ) so that some of
aluminium ions Al3 + are replaced by chromium ions Cr3 + that gives the crystal
pinkish red colour. Aluminium ions act only as host, and the actual laser light is
emitted by laser action taking place in the excited states of chromium ions. The
length of the ruby rod in laser sources may vary from 2 to 40 cm, and diameter from
0.5 to 2 cm, depending on the power output of the source. The two ends of the ruby
rod are either made flat and polished or two plane mirrors, one totally reflecting and
the other partially reflecting are put at the two ends, mirrors being held parallel to
each other and normal to the axis of the ruby rod, making the optical cavity. Often
a spring-controlled by a microscrew is attached to the fully reflecting mirror, so that
fine rotation of the microscrew may tilt in small steps the orientation of the mirror
for fine tuning of the laser wavelength.
Figure 8.20 shows the layout of a ruby laser. Optical pumping is employed in this
laser source using a xenon flash lamp, which may ether be placed near the ruby rod
or may be wound round it as a helical-shaped glass tube filled with xenon. Flashes
of light are produced when the flash lamp gas is excited by an RF oscillator powered
by a power supply. Xenon flash lamp generally emits bursts of blue-green light
(450–600 nm) of durations of few milliseconds. Parabolic mirrors are kept around
the ruby rod to focus flash light on the rod for optical pumping. Since flash light
used for optical pumping is produced in short duration pulses, the laser output of
the ruby source is not continuous but it is pulsed. Laser light is produced in pulses
of short durations of milliseconds one after the other. The energy band structure of
Cr3 + ion is shown in Fig. 8.21, where it may be seen that optical pumping using
blue-green light of xenon flash lamp may excite chromium ions from the ground
state E g to excited state bands E 3 and E 4 . Both these energy bands have very short
mean lives, respectively, ≈ 10−8 and 10−9 s and de-excite rapidly through radiation
less transition to the metastable state E 2 . The mean life of metastable band is around
10−3 s, which is much larger than the mean lives of states E 3 and E 4 . Continuous
pumping of chromium ions from the ground state to states E 3 and E 4 and the rapid
decay of these states to the metastable state E 2 , establishes population inversion
between state E 2 and the ground state E g and also between the metastable state E 2
and the excited state E 1 . Since state E 2 undergoes spontaneous photon decay mostly
to the ground state E g and only partially to state E 1 , photons emitted in spontaneous
decay initiate laser action between corresponding levels. Two groups of laser light
with wavelengths 694.3 nm and 692.7 nm are emitted (in form of bursts of laser
lights of millisecond duration) from the ruby source. The intensity of 694.3 nm laser
light is much higher than that of 692.7 nm laser. The difference in wavelengths of the
two laser groups is very small and the intensity of higher wavelength laser is much
larger; therefore, often ruby laser is treated as a single wavelength laser source.
Ruby laser is a low power and generally a pulsed laser source of red light. It
often requires a cooling system also. As such it is bulky. Initially it was used in laser
printers, etc. but now with the availability of semiconductor diode lasers which are
very small and handy, ruby laser is not in much use.
(b) Nd.YAG laser source
Neodymium YAG consists of Yttrium Aluminium Garnet (Y3 Al5 O12 ) in which some
of the Y3 + ions are substituted by Nd3 + ions. Doping level of the YAG rod is around
Fig. 8.20 Layout of a ruby laser source
Fig. 8.21 Energy band diagram of Cr3 + ion
(0.72%) that corresponds about 1.4 × 1026 Nd atoms per metre3 . In a typical Nd.YAG
source the cylindrical YAG rod length is around 10 cm and diameter around 12 mm.
General layout of the ND.YAG laser source is similar to the Cr.ruby source, except
that in this source optical pumping is done by a Krypton flash lamp, Krypton light
contains the bands of 700 nm and 800 nm which cover the absorption bands of Nd
ion.
Neodymium is a rear earth element and electron configuration of the ion is as
given below,
Nd3 + = 1s 2 2s 2 2 p 6 3s 2 3 p 6 3d 10 4s 2 4 p 6 4d 10 4 f 3 5s 2 5 p 6
Fig. 8.22 Decay scheme and laser transitions in Nd3 + ion
Electrons of partially filled level f ( = 3) are shielded by electrons of completely

filled levels 5s and 5p; however, various interactions between electrons of nearby
atoms split each level into a band. The decay scheme of Nd ion, optical pumping
and laser transitions are shown in Fig. 8.22. Optical pumping transport electrons
from the ground state band 4i9/2 to band (4f 5/2 + 2h9/2 ) which decays by radiation
less transition to the metastable band 4f 3/2 of mean life 230 µs. The metastable
band is also populated directly from the ground band by pump photons of 800–
900 nm wavelengths. Continuous pumping establishes population inversion between
the metastable band and many levels of i11/2 and i9/2 bands. Laser light of wavelengths
946 nm, 1064 nm, 1112 nm, 1116 nm, 1123 nm and 1300 nm are emitted from the
source, with different intensities. The maximum intensity is of 1064 nm transition,
while all other transitions are much weaker, therefore, the source is often referred as
a monoenergetic source of 1064 nm wavelength.
Nd.YAG source may operate both as a continuous wave (CW) source as well as
a pulsed source. A laser comprising just an active medium and resonator will emit
a pulse of laser light each time the flash lamp fires. The pulse duration will be as
long as the duration of the flash, and the peak power of the laser light will be low.
However, when a Q-switch is added to the resonator to shorten the duration of laser
pulse, output power is raised dramatically to a high value. A Q-switch allows the
build-up of population inversion in the active medium by retarding the laser light
emission, which results in short pulse of laser light of very high power. The process of
building-up of population inversion may be compared to the charging of a capacitor
to a high value of charge (potential) before its discharge. In Q-switching mode, output
power of 250 MW for pulse duration of 10–25 ns may be achieved. In pulse mode
operation the source requires either forced air or water jet cooling.
Nd ion may be doped in other host materials also, like in yttrium lithium fluo-
ride (YLF) which emit laser light of 1047 and 1053 nm wavelengths; yttrium
orthovanadate (YVO4) that produces laser light of 1064 nm and glass.
Nd.YAG laser is one of the most used lasers; it is used in ophthalmology to
correct vision disorders, in ablation of malignant and benign lesion in different body
parts, and in cosmetic surgery, etc. It also has wide applications in industry, live
engraving, etching, laser cold peening, cutting and drilling holes in metal sheets,
in nonconventional rapid prototyping process called laser engineered net shaping
(LENS), etc.
(ii) Solid state semiconductor diode laser source
Semiconductor diode lasers, also called quantum well laser, are not much different
from a light-emitting diode (LED). The two important differences between the LED
and diode laser are: (i) diode lasers use some direct semiconductor, like GaAs and
have heavily doped p- and n-sides, while in LED the doping levels are not so high.
(ii) Laser diode operates at high forward current, larger than a threshold value, while
LED operates at lower forward currents. Emission of monoenergetic photons both
in LED and the laser diode is through the process of electron–hole recombination;
however, the process of recombination is spontaneous in LED while it is simulated
emission in case of laser diode.
Laser diodes are fabricated using direct semiconductor material like, GaAs and
doping the n- and p-sides heavily so that they are represented as n+ and p+ . There can
be two types of laser diodes: (a) homojunction, where the material of the n-p junction
is same and (b) heterojunction, where two different semiconductor materials are
used to fabricate the pn junction. However, the working of the two types of laser
diodes is similar. Let us first study the working of a homojunction laser diode. A
homojunction laser diode is made by heavily doping the two sides of a single crystal
of a direct semiconductor, like GaAs, with p- and n-type impurities with a very thin
(≈ 1 µm) depletion layer.
Figure 8.23a shows the conduction and valence bands of a heavily doped pn
junction when no external bias is applied to the junction. In case of a pn junction where
both sides are heavily doped, Fermi level passes through the conduction band on n-
side and through the valence band on the p-side, thus developing a potential barrier
across the junction. As a result of the potential barrier majority carrier electrons from
the n-side and holes from the p-side could not cross from one side to the other. Forward
bias voltage, when applied across the junction, opposes the inbuilt potential barrier
across the depletion layer and pushes majority charge carriers across the junction (see
Fig. 8.23b). Electrons and holes crossing the thin depletion layer recombine emitting
monochromatic light in case of ordinary LED. In case of LED recombination process
is spontaneous and light photons are emitted in different directions with random
phases. However, on increasing forward current to a higher value, larger than the
threshold current, large number of majority carriers from both, the n- and the
p-sides get accumulated in a small region around the physical junction, creating
population inversion. A photon of spontaneous recombination may work as the seed
photon to initiate laser action. Later, laser photons initiate subsequent stimulated
photon emissions, providing a gain factor to the lasing medium, which in this case
is the depletion layer. Almost two-dimensional depletion layer also serves as optical
cavity to develop sustained laser oscillations and amplify the laser yield. However,
no external mirrors to increase the path length for sustained laser oscillations are
required in this case. Out of the six sides, four sides of the depletion layer (the
optical cavity) are exposed to the air of refractive index 1, while the refractive index
of GaAS for laser light is around 3.8. Hence the reflectivity R of the depletion layer–
air interface is R = (3.6−1)
2
(3.6+1)2
= 0.3 which is reasonable for laser light to get reflected
back and forth between the two edges of the depletion layer cavity. Out of the four
air-exposed faces of the depletion layer, two opposite faces are cleaved and their
parallelism is ensured. The other two faces are left rough. Laser oscillations build
up between the two parallel cleaved faces and are taken out from one of them as
indicated in Fig. 8.24.
When forward current is low, gain of resonance optical cavity gets compensated
by the losses and, therefore, the overall medium gain remains below the critical
gain value. Increasing forward current above the value called threshold current, the
medium gain becomes larger than the critical gain giving rise to the build-up of laser
oscillations in the cavity.
Fig. 8.23 Conduction and valence bands of a heavily doped, a unbiased pn junction, b forward
bias pn junction
Fig. 8.24 Layout of a laser

diode
There are two main drawbacks of a homojunction diode laser;

(i) since the cross section of the depletion layer ‘d’ is very small (≈ 1 µm), of the
order of the wavelength λ of light, the divergence angle θ of laser beam coming
out of the diode is large on account of diffraction, θ = λ/d.
(ii) The efficiency of a homojunction laser diode is relatively small. It is because
for high efficiency it is required that (a) both the injected charged carriers and
laser photons emitted by recombination remain confined within the junction
region (depletion layer) in order to initiate repeated stimulation emissions and
(b) electrons/holes should travel small and equal distances for recombination.
Though condition (a) in homojunction lasers is achieved by employing large
density of forward current but condition (b) is not accomplished in homojunction
diodes since electrons travel different distances before recombination.
Some advantages of diode laser are: they are small, cheap, can be produced in
large number and are easily scalable. Laser light of different wavelengths may be
obtained from laser diodes by choosing appropriate material.
Heterojunction diode laser A homojunction and two types of heterojunction diode
lasers are shown in Fig. 8.25. A double heterojunction semiconductor laser is made by
sandwiching a thin layer of GaAs between two layers of some ternary semiconductor
material with general chemical structure Ga1−x Alx As (here x is the mole fraction of
Al and 1 − x the mol fraction of Ga). The ternary compound must have a refractive
index (for laser light) smaller than the refractive index of GaAs for the same light
and the energy band gap larger than that of GaAs.
The lower refractive index of surrounding ternary compound reflects back the
laser light photons to the active region by the process of total internal reflection, thus
confining photons within the active GaAs region. The larger band energy gap pushes
back electrons to the active region of GaAs. In this way, both photon and charge carrier
confinement along with nearly same travel distance of electrons before recombination
are achieved in case of heterojunction laser diode. Heterojunction diode lasers are
much more efficient that means that in them population inversion is achieved at much
lower density of forward current.
Fig. 8.25 Homo and two types of heterojunctions
8.7.2 Dye (Liquid) Laser Source
Dye lasers are typical example of liquid state lasers. An organic dye is used
as the active medium in a dye laser. Frequently used organic dyes for liquid
lasers are Rodamine 6G, (Xanthene), Anthracenes, Oxazines, Coumarin, DCM (4-
dicyanomethylene-2-methyl-6-pdimethaylaminostyryl-4H pyran), etc. These dyes
are dissolved in appropriate solvent, and the solution acts as the gain medium. The
dye solution is optically pumped and the dye molecule absorbs a band of wave-
lengths of average energy E, from the incident light of the flash lamp. The excited
dye molecules revert back to the ground state, emitting a band, called fluorescence
band, of average energy E 1 which is less than the energy E of absorption band. The
difference of energy (E − E 1 ) is generally dissipated in non-radiative collisions. Dye
molecules are generally long-chain big molecules, and only a part of the molecule
takes part in excitation and de-excitation.
Energy band structure of liquid dyes is classified in terms of singlet S and triplet T
bands. S-bands that correspond to angular momentum state = 0, and triplet bands
corresponding to angular momentum state = 1; have stacks of bands of increasing
energy, separated from each other with certain energy gaps, as shown in Fig. 8.26.
Each singlet and triplet band has a cluster of closely spaced vibrational and rotational
levels.
A general property of s-bands is that higher levels in a given s-band decay to the
lowest energy level of the band by non-radiative de-excitation in a very short time of
the order of 10−11 s. For example if levels of band s1 are excited by optical pumping
(by shining light of appropriate frequency), all levels from b to B will get excited but
within a time span of 10−11 s all higher levels up to B will revert back to the lowest
Fig. 8.26 Typical energy level structure of a dye molecule
level b. It is also observed that the lowest level of each s-band is relatively longer
lived as compared to the higher levels. In case of band s1 of the figure, the mean life
of the lowest level ‘b’ is around 10−5 s, an order of 106 longer than the mean lives of
higher levels of the bands. In a way, the lowest level of a s-band may be considered
like a metastable state (though it is not a metastable state since its life time is still
very small of the order of 10−5 s).
In normal case, when the dye solution is irradiated with a source of light of
appropriate energy, molecules in ground state of band s0 absorb energy from the
incident light and shift to various levels of band s1 . Excited molecules in levels above
‘b’ quickly lose extra energy by collisions, etc. and de-excite by non-radiative mode
to reach level ‘b’ in 10−11 s. Once most of the excited molecules get accumulated in
level ‘b’, they decay by emitting a bunch of photons to level A, at the head of band s0
in around 10−5 s. This emission of photon bunch from b to A is called fluorescence.
Fluorescence photons are incoherent, and their mean frequency is less than the mean
frequency of the photons absorbed from the incident light. That means that if the
system is irradiated, say by blue light then fluorescence light may be of red colour.
In fluorescence emission time delay between absorption and re-emission of light is
very small, generally less than a few microseconds, however, if the time delay is
measurable, say of the order of few seconds or larger, the process of light emission
is termed as phosphorescence.
Continuous optical pumping by appropriate flash lamp increases the population of
excited molecules on level ‘b’ at the cost of ground state level ‘a’ leading to population
inversion. A fluorescence photon may work as seed to initiate laser emission in
the early stages and the triggering is taken over by laser photons later on. In dye
laser, laser photons are made to circulate in a resonance cavity made by putting two
parallel mirrors, one fully reflecting and other partially reflecting, at the two ends
of the tube containing the solution of the dye. Dye solution works as gain/active
medium. If the medium gain is larger than the critical gain, laser oscillations build up
in the cavity and strong laser beam emerges from the partially reflecting mirror. Dye
laser source consists of a quarts or glass tube (≈ 10 cm, ≈ 0.5 cm diameter) with
dye solution. The solution tube is converted into a resonance cavity by putting two
parallel mirrors at the two ends of the tube. The cavity tube is placed at one focus
of an elliptical reflecting envelop at the other focus of which a flash lamp is placed.
Elliptical envelop concentrates the light emitted by the flash lamp on the cavity tube.
Flash light is operated by the usual electronic circuits that generate flashes of light of
few milliseconds each. Laser light beam emerges from the partially polished mirror.
The frequency (or wavelength) spread of dye laser is considerably larger than solid
state lasers. This results in the production of large number of longitudinal modes
in cavity. Any of these modes may be selected by frequency selection components
associated with the cavity. Dye lasers are, therefore, tuneable to different frequencies.
This is the big advantage of a dye laser.
Dye lasers are always used in pulsed mode. The reason for this is the energy gap
(E 1 ) between the head of T1 and the bottom of T2 bands. In most dyes, this energy
difference is of the same order as the energy of photons emitted in laser beam. If there
are sizable numbers of excited molecules in state T1 , then large number of emitted
laser photons may be absorbed by molecules in state T2 and output laser flux may
die out. As may be observed in Fig. 8.26, T1 band is fed by the decay of level ‘b’
with a mean life of around 5 s, i.e. the feeding transition is very slow. It means that
if laser is produced in pulses of short durations, so that population of level ‘b’ is not
allowed to grow for long time, there will never be enough molecules in state T1 and
the chance of absorption of laser photons will be eliminated. This is the reason why
dye lasers are always used in pulsed mode.
The output wavelength of dye lasers generally varies from 390 to 1000 nm and
output power may be between milli-Watt to 1 W. Typical beam diameter may be
0.5 mm with beam divergence from 0.8 to 2 mrad.
The main advantage of dye laser is that laser beams of any wavelength ranging
from infrared, visible, and up to near ultraviolet may be produced using a dye laser,
that is why they are called tuneable lasers. Dye lasers have high efficiency ≈ 25%,
small divergence, small beam spot and high power. However, they are always pulsed
source.
8.7.3 Gas Laser Sources
Gas lasers have some gas as the active medium. Gas lasers, depending on what specie
produces lasing action, may be divided into three types: (i) atomic, (ii) molecular
and (iii) ionic.
(i) Atomic gas laser
Helium–neon laser is the best example of atomic gas laser. It was the first continuous
wave (CW) laser built by Ali Javan, Bennett and Horriot at Bell Telephone Lab in
1961. In He–Ne laser source a mixture of He and Ne gases in the ratio 10:1 is taken
in a small glass or quartz tube of approximately 10–100 cm in length and 5–10 mm

in diameter at about 10 Torr of pressure. An electric discharge is then passed through
the gas mixture that dissociates gas molecules into atoms and ionises some atoms.
Ionisation electrons get accelerated in the electric field of the applied high voltage
and collide with gas molecules.
Collision of helium atoms with accelerated electrons in the discharge tube excites
them to 2S1 and 2S3 levels as shown in Fig. 8.27. Since the number of Helium atoms
is much larger than that of neon atoms, large numbers of excited helium atoms collide
with ground state atoms of neon and transfer their excitation energy to neon atoms
that get excited to 3S2 and 2S2 levels which are very close to the excited levels of
Helium. Once transferred their excitation energy to neon atoms, helium atoms revert
back to their ground state, collide again with accelerated electrons, get excited again,
transfer excitation energy to neon atoms, revert back to ground state and this cyclic
process goes on. The excited states of 3S2 and 2S2 of neon have longer mean lives
in comparison to other excited levels of helium and, therefore, behave as metastable
states. Continuous transfer of energy by excited helium atoms to ground state neon
atoms establishes population inversion between neon energy levels 3S2 and 3P4 ; 3S2
and 2P4 and between 2S2 and 2P4 . Spontaneous decay photons between the above
mentioned levels work as seed photons to trigger laser action. The stimulation photon
emission is then triggered by laser photons that circulate in the optical cavity created
by putting two parallel mirrors, one partially polished, at the two ends of the He–
Ne filled discharge tube. Helium–neon gas mixture serves also as the gain medium.
Three laser groups with mean wavelengths 3391.2 nm; 1152.3 nm and 632.8 nm are
emitted by the He–Ne laser source.
Fig. 8.27 Energy level diagram of helium and neon

He–Ne laser is very frequently used in laboratories as demonstration laser as

it is cost effective, rugged, has no problem in continuous running for hours, heat
dissipation being no issue but deliver only small power output of the order of
milliwatts.
(ii) Carbon Dioxide Molecular gas laser
CO2 molecular laser is one of the most frequently used continuous wave high power
laser that has large number of industrial applications. CO2 laser generally operates
in the infrared region of electromagnetic spectrum which is invisible to human eye
with wavelength near 10−6 m and power as large as 100 W/m2 .
Energy level spectrum of molecules is more complex as compared to atoms or
ions. In molecules there are electronic excitations, vibrational excitations and rota-
tional excitations, in decreasing order of energy. Population inversion in CO2 laser
is achieved between vibrational levels. CO2 molecule may have three types of vibra-
tional modes, namely (i) symmetric stretching denoted by i00, where i stands for
an integer value of energy, (ii) bending mode, denoted by 0j0, integer j denoting
the energy and (iii) antisymmetric stretching, denoted by 00k, integer k denoting the
value of energy. These vibration modes for CO2 molecule are shown in Fig. 8.28a.
Nitrogen molecule having two identical atoms does not have bending mode. Energy
level diagram for vibrational motion of N2 and CO2 molecules is shown in Fig. 8.28b.
Carbon dioxide laser sources may be of two types; sealed tube type and continuous
flow type. In sealed tube-type source a mixture of N2 and CO2 gases each around
10–20% by volume and about 80–90% of Helium gas, at about 10 Torr pressure are
taken in a glass discharge tube of few millimetre diameter and of length from few
Fig. 8.28 Vibrational energy levels of N2 and CO2 molecules

ten of cm to 1 m. Two electrodes, one at each end of the discharge tube, are used to
produce gas discharge in the tube. Since CO2 laser source may deliver large power,
it also consumes lots of electrical energy, producing considerable amount of heat.
Both N2 and CO2 gases in discharge tube take part in lasing action, while helium
simply helps in heat dissipation from the interior of the discharge tube to its walls
and to some extent in de-excitation of lower vibrational levels (020 and 010) of CO2
molecules. Closed discharge tube is wrapped with coils through which chilled water
is circulated to remove heat from the tube walls. In gas flow-type source, mixture
of N2 and CO2 gases is made to flow through the discharge tube at a constant rate,
outgoing gases takeaway heat produced in the discharge tube and hence no helium
is required in the gas flow type source.
The gas mixture in discharge tube also serves as the gain medium, and the
discharge tube is converted into a resonance cavity by putting two mirrors, one
partially polished, parallel to each other and normal to the tube axis. Since glass
absorbs laser light emitted by CO2 molecule, reflecting mirrors and other equipment
used inside the discharge tube are made from materials like Ge, GaAs, ZnS, etc.
which are transparent to laser radiations.
Working of the CO2 laser source may be easily followed by looking at Fig 8.28b.
Once discharge is produce in the tube by applying large voltage between the two elec-
trodes, there is large number of free electrons accelerated under the high electric field
of applied voltage. Energy of electrons in discharge tube is much larger than 0.3 eV,
which is the energy of the first vibrational excitation level of N2 molecule. Collisions
between energetic electrons and N2 molecules excite large number of molecules to
the first excited state. The first vibrational excited state of N2 is metastable; there-
fore, excited N2 molecules live long enough to collide with ground state molecules of
CO2 and transfer their excitation energy to the CO2 molecules. CO2 molecules also
have vibrational excited state (003) (asymmetric vibration mode at energy 0.3 eV)
very near to the in energy to the first vibrational state of N2 . Thus collisions of N2
molecules with accelerated electrons excite them to their first vibrational level and
collisions between excited N2 molecules with ground state molecules of CO2 , excite
CO2 molecules to their vibrational level 003. There are two vibrational levels 200
and 020 of CO2 at 0.2 eV energy, which are almost empty. Continuous pumping (by
N2 molecules) of CO2 molecules from ground state to 003 level establishes popu-
lation inversion between levels 003 and 200 and between levels 003 and 020. Since
molecules in level 200 lose their energy by collision with molecules in level 020, and
level 020 readily depletes through diffusion and in closed tube source in collision
with Helium molecules, the state of population inversion is maintained. CO2 laser
source produces two laser lights of wavelengths 9600 nm (9.6 µm) and 10,600 nm
(10.6 µm).
The main advantages of CO2 laser are simple in construction, both CW and pulsed
outputs, very high power output, high efficiency.
CO2 laser is extensively used in treating health problems related to gynaecology,
genitourinary, dental, orthopaedic, hepatic and cardiovascular surgery. It is consid-
ered to be the mainstay of laser neurosurgery. They are used for cutting, dissection
and coagulation of wide range of tissues.
High power CO2 laser finds application in industry, in material processing,

welding, drilling, cutting, soldering, etc.
Because of its low absorption in atmosphere, powerful beams of CO2 laser are
used for remote sensing and for open air communications.
(iii) Argon ion laser
An argon ion laser source is made by filling natural argon gas in a tube of some ceramic
material like beryllium oxide, that is opaque to laser radiations and can withstand
high temperatures. The tube is fitted with two hollow electrodes to produce plasma
discharge. The density of Ar+ ions in the plasma is high, and to further increase
it, a solenoid magnetic coil is wrapped around the plasma tube which confines the
plasma and fast-moving electrons in the central part along the axis of the tube. A
typical source may have plasma tube up to 1 m long that may generate laser with
output power of 2–5 W consuming several tens of kilowatt input power. The efficiency
of the source is very poor, and chilled water cooling is required to dissipate the heat
developed in the source. The plasma tube is fitted with Brewster windows at the
two ends, as ordinary glass absorbs the plasma light and also high temperature is
produced around the tube. Two parallel mirrors, one partially transparent are put
after windows normal to the plasma tube axis, as shown in Fig. 8.29. A prism is
inserted between Brewster window and the totally reflecting mirror to tune laser rays
of different wavelengths.
The level scheme of excited states of Ar ion and pumping of excited-ion levels by
collision of high-energy electrons with argon ions is shown in Fig. 8.30. Continuous
pumping of excited states by electron collisions establishes population inversion
between several excited states, resulting in the emission of groups of lasers. Some
prominent laser groups are shown in Fig. 8.30. The source may produce laser beams
of as many as 35 different frequencies; however, the most prominent laser beam has
wavelength of 514.5 nm.
Fig. 8.29 Layout of an argon ion laser source. Magnetic solenoid and cooling system are not shown
in the figure
Fig. 8.30 Level scheme for argon ion laser system
Argon ion laser is frequently used in forensic medication, holography, Raman

spectroscopy, high-speed printing, in laser entertainment shows, etc.
8.7.4 Excimer Laser
There are some molecules, like ArF, KrF, XeCl, etc. that are stable in their first excited
states but dissociate in their ground states. Such molecules are called excimers.
Population inversion in such molecules is easy to obtain as in natural way the number
of molecules in ground state are negligible as compared to those in the first excited
state.
A typical excimer laser source may be made by filling an inert gas like argon and
some halide like F2 in a discharge tube fitted with electrodes. High voltage applied
across the two electrodes set in discharge in the tube producing Ar+ positive and
F− negative ions which combine forming ArF* molecule in its first excited state.
Excited ArF* molecules revert to the dissociative ground state initiating laser action.
The efficiency of excimer lasers is of the order of 20% and they emit laser radiations
in the wavelength range of 120–500 nm with peak power of around 200 W.
8.7.5 Mode Locking
A typical laser source consists of an optical resonating cavity formed by to coaxial

mirrors and a gain medium. The wavelength band over which laser oscillations may
occur is determined by the wavelength region over which the gain of active medium
exceeds resonator losses. Mostly the above condition is fulfilled for several longi-
tudinal (or axial) modes of closely spaced wavelengths/frequencies which produce
stationary oscillations in the resonator cavity. The laser output in such cases consists
of laser radiations of closely spaced wavelengths/frequencies with decreasing ampli-
tudes (see Fig. 8.12c). Let us assume that in a typical case there are N axial modes
that produce sustained oscillations in a resonator of length L. Now all the N-modes,
in general, will have different amplitudes An , angular frequencies ωn and phases δn ,
where n varies from 0 to N. Since all the three parameters are functions of time,
therefore, in general modes will be incoherent. The total output (amplitude) of such
a laser will be a linear combination of different modes and will be given by
∑
N
A(t) = (A)n ei(ωn t+δn ) (8.14)
0
If there is nothing that fixes the three parameters, amplitudes, frequency and
relative phases of the N different modes, then the output amplitude A(t) will vary in
an uncontrolled way. However, if the different modes are forced to maintain equal
frequency spacings with a fixed phase relationship to each other, the output with
time will vary in a well-defined manner. The laser is then said to be mode-locked
or phased-locked. The form of the output will depend on which axial modes are
oscillating and what phase relationship is maintained. It is possible to obtain: an
FM modulated output, a continuous pulse train, a spacially scanning laser beam or
a ‘machine gun’ output where pulses of laser light appear periodically at different
spacial positions on the laser output mirror. Mode locking is often used to produce
short-duration, high-intensity bursts of laser radiations at a given repetition rate.
There are several ways of doing mode locking; however, details of that are beyond
the scope of the present discussion.
8.7.6 Q-Switching
Q-switching, also called ‘giant pulse formation’ is a technique of producing very

high-powered, peak power of the order of gaga watts, laser pulses of very short
durations, of the order of picoseconds, from a normal laser source. In comparison
with mode-locking technique, which is also a method of producing laser pulses, the
peak power of the output laser pulse in Q-switching is much higher, with lower pulse
repetition rate, and longer pulse durations. In contrast with mode locking which
provides a train of very short duration laser pulses, Q-switching provides a single
laser pulse of very high peak power and of short duration.
Q-switching is achieved by putting an attenuator in resonator cavity. The attenu-
ator inhabits the round trip motion of laser photons between the two resonator mirrors.
As a result, the gain of the active medium reduces and the resonator quality factor
(also represented by Q) decreases. Attenuator converts a high-quality resonator into
a low-quality resonator. However, the number of excited atoms in active medium

keeps on increasing on account of continuous pumping. Thus the energy contents of
the active medium keep increasing during the attenuator action, reaching a saturation
value. At this point the medium is said to be gain saturated. Once the gain satu-
ration is achieved, the attenuator stops restricting the passage of laser photons and
allow them to undertake round trips between the two mirrors initiating stimulated
emission of laser radiations. Since large amount of energy is already stored in the
active medium (in the form of excited atoms produced by pumping), the intensity
of the laser radiations in resonator builds up very quickly. Quick build-up of laser
intensity in resonator causes the energy stored in the medium to be depleted almost
as quickly. Net result is a short laser output pulse, called the giant pulse that may
have very high peak intensity/power. Attenuators are generally electro-optical, or
opto-acoustic devices. Saturable light absorbers may also be used for attenuation of
medium gain.
Some characteristic properties of important lasers are listed in Table 8.1.
Table 8.1 Properties of some important laser sources

Active Wavelength Power/ Beam Divergence Efficiency Cooling
medium λ (nm) energy diameter (mrad) (%)
(nm)
Gain medium solid state
Nd.YAG 1064 Up to 10 kW 0.7–10.0 0.3–25 0.1–5 Air/water
Nd.Glass 1060 < 100 W 3–25 3–10 1–5 Air/water
Ti.Sapphire 660–1000 2W Few mm Few mrad 0.1–5 Air/water
Er.fiber 1550 1–100 W < 10 Few mrad 1–5 Air
GaAs, 780–900 1 mW–10 W Diverges 200 × 600 1–50 Air/heat
GaAlAs rapidly oval-shaped sink
InGaAsP 1100–1600 1 mW–1 W Diverges Oval shaped 1–20 Heat sink
rapidly
Gain medium gas
He–Ne 632.8 0.1–50 mW 0.5–2.5 0.5–3 < 0.1 Air
CO2 10,600 3 W–20 kW 3–50 1–3 5–15 Gas flow
Gain medium gas excimer
Argon 193 Up to 50 W 4 × 2 to 25 2–6 <1 Air/water
fluoride × 30
rectangular
Gain medium liquid
Different Tuneable 20 mW to 1 1–20 0.3–2 1–20 Water/dye
dyes 300–1000 W flow
8.8 Some Applications of Lasers 425
8.8 Some Applications of Lasers
It is difficult to give an exhaustive list all laser applications; however, some important
laser applications are discussed here.
(a) Laser holography
A conventional photograph carries two-dimensional information of the light inten-

sity (that is proportional to the square of the light amplitude) scattered by a three
dimensional object. All information about phase relationships of scattered light, that
tells about the depth, the third dimension, is lost in ordinary photography. Holog-
raphy is the technique of taking photographic records where both intensity and phase
attributes of scattered light are recorded. A holograph also gives some feeling about
the third dimension of the object.
Hologram is a recording in a two-dimensional or three-dimensional medium of
the interference pattern formed when monochromatic light from a point source inter-
ferers with the monochromatic light of same wavelength scattered by an object. The
interference pattern or hologram when illuminated by the same monochromatic light
produces the diffraction pattern forming an image of the object that is indistinguish-
able from the object. Since laser lights are highly monochromatic and laser sources
are very nearly point sources, lasers are very often used for making holograms.
Holograms may be of three different types; (a) reflection hologram, (b) transmission
hologram and (c) hybrid hologram.
As already mentioned, making of a hologram requires two beams of same
monochromatic light coming from the same point source. The beam of light that
comes directly from the point (laser) source and illuminates the recording photo-
graphic emulsion is called the reference beam, and the other beam reaching the
photographic emulsion after scattering from the object is called the object beam.
Figure 8.31 shows the basic layout of a system for recording of a hologram. As
shown in the figure, the parallel beam of monochromatic light from a point laser
source is split into two parts, one part on reflection by the mirror directly illuminate
the recording photographic emulsion and forms the reference beam. The other part
of the laser beam illuminates the object and laser light scattered by the object (called
object beam) interferes with the reference beam at photographic emulsion, producing
a permanent record of interference pattern on photographic emulsion.
Reconstruction of the object image from the hologram can be done as shown in
Fig. 8.32. The same laser source that was used for making the hologram is used to
illuminate the hologram at the same angle at which the reference beam hit it while
recording the hologram. The incident laser beam gets diffracted from the interference
pattern recorded on hologram and the secondary wavelets emitted from interference
pattern undergo constructive and destructive interference to make a real and a virtual
images of the object. Superimposition of images produces three-dimensional image
of the object with all details. In a hologram light scattered by all parts of an object
produces the interference pattern at each point of the hologram. As such each small
part of hologram contains information about whole of the object.
Fig. 8.31 Recording of a hologram
Fig. 8.32 Reconstruction of image from hologram
One major advantage of hologram is that each fragment or a small piece of a

hologram can reproduce the complete image of the object while a piece of an ordinary
photograph cannot reproduce the complete image of an object.
(b) Writing and reading of digital data on compact disc (CD) using laser beam
A master compact disk (CD) is made of three components; plastic on which digital
data is written using a laser beam, reflecting aluminium layer and a protective polycar-
bonate plastic as shown in Fig. 8.33. The master disc is ‘burned’ with a laser beam that
etches bumps, called pits, on the plastic surface. A bump represents binary number
Fig. 8.33 Basic structure of

a compact disc
0. As such when the laser beam etches a bump a zero is stored on the master disc.
The flat unburned (or unetched) plastic surface represents binary number 1. The flat
unetched area representing 1 is called land. Thus desired digital data is stored on the
master disc in the form of pits and lands. Once the master disc is made it is used to
stamp millions of plastic duplicates—the CD’s sold in the market. After stamping
each disc is coated with a reflecting aluminium layer and covered with a protective
polycarbonate layer.
Laser light is used to read the data from a CD. A thin laser beam scans the CD
surface and gets reflected from the pits and lands covered with reflecting aluminium.
The intensity of the back reflected laser light is different for reflection from a pit and
a land. A photodiode records these variations in the intensity of back reflected laser
light and the data on intensity variation is converted into 0 and 1.
(c) Military/defence/armament applications of laser
Military uses of lasers include target designation and ranging, defensive counter
measures, directed energy weapons and communication, etc.
A low power laser pointer is used to indicate a target for a precision-guided
munition, generally lunched from an aircraft. The guided munition adjusts its flight-
path to home on the laser light reflected from the target, enabling a great precision
in aiming. The target designator laser beam is pulsed at a typical rate that is sensed
by the munition to avoid any confusion with other laser beams present in the area.
Powerful beams of laser light (more than 100-kW power) have been used to
destroy approaching ballistic missiles. Laser beams have also been used to destroy
land mines and unexploded explosives scattered in the war zone. Laser beams have
been used for counterdefensive measures; low-power infrared countermeasures use
lasers to confuse the heat seeker anti-aircraft missiles while high power boost-phase
interceptors use lasers to find, track and destroy ballistic missiles. Some weapons
simply use a laser to disorient enemy army personals. Lasers have also been used
as a tool to enhance the targeting of other weapon systems, for example, a machine
gun or a rifle barrel is fitted with a laser torch that emit a fine beam of visible laser
light parallel to the barrel that may travel up to few kilometres. The laser light spot
may be aligned with the target for better aim, taking into account wind direction and
speed and trajectory of the fired bullet/shell.
(d) Industrial and commercial applications
Depending on the power of the laser source, industrial applications of laser may be
divided into two types; material processing and micromaterial processing. In material
processing laser sources of more than 1 kW power are used. Laser sources in the
range of 100–300 W are primarily used for pumping, plastic welding and soldering
applications. In applications like, metallic sheet cutting, brazing, metal welding etc.
laser sources with power larger than 300 W are employed. Multiple kilowatt lasers are
used for hardening, deep penetrating welding, cladding etc. Lasers are also used for
micromaterial processing, like fabricating screens for smart phones, table computers,
and LED TVs.
Lasers are used for optical communication over optical fiber or in free space, in
guidance systems like laser gyroscopes, barcode readers, laser engraving of printing
plates, writing subtitles on motion picture films, in consumer and industrial imaging
instruments, etc.
(e) Medical applications
Lasers are now extensively used for cosmetic surgery, like removing of tattoos, scars,
stretch marks, sunspots, wrinkles, birth marks and unwanted hairs. Ruby (694 nm),
Nd.YAG (1064 nm), alexandrite (755 nm) and pulsed diode array (810 nm) lasers are
used for dermatology applications. Soft tissue surgery is done mostly using CO2 and
Er.YAG lasers. Laser scalpels are now used for general surgery and in gynaecological,
urology and laparoscopic surgeries. Laser knifes are used for no-touch surgery of
brain and spine tumours. In dentistry, lasers are used for caries removal, in endodontic
and periodontal procedures etc.
Solved Examples
SE8.1 A resonant cavity in some laser source is made by putting two parallel
mirrors 1 m apart. Calculate the frequency separation Δν between different
axial modes. Will frequency separation be different for lasers of different
frequencies?
Solution Frequency difference Δν between successive axial modes is given as,
c/μ
Δν =
2L
where c is the velocity of light, μ the refractive index of active medium and L is the
length of the cavity. Since the value of μ is not given in the problem, we take it 1,
therefore,
3 × 108 m/s
Δν = = 1.5 × 108 cycles per second = 150 MHz
2 × 1m
Frequency separation does not depend on the frequency of the laser; it depends
only on the length of the resonant cavity.
SE8.2 A typical laser source emits laser light of 700.0 nm wavelength with spread
of 0.50 nm. Calculate the possible number of axial modes in a cavity of 5 cm
length filled with gain medium of 1.5 refractive index.
Solution Let us first calculate the frequency uncertainty of the laser light that has
wavelength uncertainty of 0.50 nm. The frequency spread δν corresponding to
wavelength spread δλ = 0.5 nm is given by,
c 3 × 108
δν = δλ = 0.5 × 10−9 = 309 × 109 Hz.
λ2 700 × 10−9
2
Next we calculate the frequency separation Δν between successive axial modes,
c/μ 3 × 108 ms
Δν = = = 4 × 109 Hz
L 1.5 × 5 × 10−2 m
δν 309×109
Number of axial modes is = Δν
= 4×109
= 77.5 modes ≈ 77 modes.
SE8.3 A laser beam of 50 mW power having a circular spot of radius 1.0 mm is
allowed to fall on a totally reflecting mirror for 10 h. Calculate the intensity
of the laser light, the linear momentum imparted by the beam to the mirror,
average force exerted by the beam and the total electrical power consumed
by the source if the efficiency of the source is 10%.
Solution The power of the beam W = 50 mW = 50 × 10−3 J/s. This power is
2
contained in a circular beam spot of area A = πr 2 = π × 1 × 10−3 . Intensity I
of the beam is;
W 50 × 10−3
I = = 2
J s−1 m−2 = 15.91 × 103 J s−1 m−2
A π × 1 × 10−3
Each light photon of energy E = hν carries a linear momentum hν/c = E/c. If

there are N photons each of frequency ν, then total linear momentum carried by N
photons L = Nhν/c = total photon energy/c. As given in the problem 50 × 10–3 J/
s is the energy deposited by photons on the mirror, hence in 10 h amount of energy
deposited on mirror is given by,
E T = 50 × 10−3 × 10 × 60 × 60 = 572 J
Total linear momentum imparted by light photons in 10 h L imp = EcT .

Since mirror is totally reflecting all photons after hitting the mirror get reflected
and takes away the same total momentum but in opposite direction, hence change in
momentum in 10 h

ET ET 2E T 2 × 572
ΔL = − − = = = 381.3 × 10−8 kg m s−1
c c c 3 × 108
Total linear momentum imparted to the mirror by laser beam in 10 h = 381.3 ×

10−8 kg m s−1 .
Average force entered by the laser beam on mirror Fav =
Rate o f change o f momentum.
or
ΔL 381.3 × 10−8
Fav = = = 10.59 × 10−11 N
10 × 60 × 60 10 × 60 × 60
The efficiency of the laser source is given as 10%, therefore,
output light power 50 × 10−3 W

10 = × 100 = × 100
input electrical power Input electric power
Input electric power = 50 × 10−3 W × 10 = 500 mW.

SE8.4 A laser beam of wavelength 7000 Å and of aperture 0.05 cm is flashed on a
target 5 km away. Calculate the angular and area spread of the beam when
it reaches the target.
Solution Given that λ = 7000 Å = 7000 × 10−10 m = 7 × 10−7 m
and aperture ‘d’ = 0.5 cm = 0.05 × 10−2 m = 5 × 10−4 m.
Distance D of the target = 5 km = 5 × 103 m.
−7
Angular spread dθ = dλ = 7×10
5×10−4
= 1.4 × 10−3 rad.
2
Area of spread A = (Ddθ )2 = 5 × 103 × 1.4 × 10−3 = 49 m2 .
SE8.5 A laser beam of wavelength 700 nm is chopped in pulses of 1 ns using a
shutter. Calculate the line width, the bandwidth and the coherence length of
the pulsed laser.
Solution Coherence length = c.pulse width = 3 × 108 × 1 × 10−9 = 3 × 10−1 m =
0.3 m.
Band width Δν = pulse1width = 1×101
−9 = 1 × 10 Hz.
9
Line width Δλ = λ Δν = (
700×10 )
−9 2
× 1 × 109 = 18.3 × 10−13 m.
2
c 3×108
Problems
P8.1 A laser beam of wavelength 700 nm and power 1 mW was focused on an spot
of area 40 × 10−14 m2 . What is the intensity of the focused beam?
ANS 2.5 × 109 W/m2
P8.2 In case of a ruby laser the laser beam diameters at 2 m and 4 m from the source
were 2 mm and 3 mm. What is the angle of divergence of the beam?
ANS 2.5 × 10−4 rad
P8.3 A laser beam of wavelength 800 nm and aperture 0.5 cm is sent to an object
4 × 108 m away. Calculate angular spread and area spread of the beam at the
distant object.
ANS 16 × 10−5 rad; 409.6 × 1010 m2
P8.4 The optical cavity of a laser source has the active medium of refractive index
1.75 and length 5 cm. The central wavelength of the laser resonating in the
cavity is 700 nm with spread 0.5 nm. How many axial modes will be resonating
in the cavity?
ANS 179 axial modes
SA8.1 What is meant by population inversion? Is population inversion a state

of thermal equilibrium? If not then how it could be maintained in a laser
source?
SA8.2 In a laser source the upper lasing level is mostly a metastable state, Why is
it necessary?
SA8.3 Suppose a laser source does not have a resonant cavity, what will happen
to the source, will lasing action take place in it? Will it deliver an intense
beam of laser light?
SA8.4 The active medium of a resonant cavity has a gain factor. In what form does
this gain remain distributed in the active medium?
SA8.5 List the possible reasons for the spread of wavelength/frequency of a laser
source.
SA8.6 On what factors the frequency separations between successive axial modes
in a resonant cavity depend?
SA8.7 List important characteristics of laser light explaining two of them.
SA8.8 Intensity of an image formed using laser source is very high, explain why?
SA8.9 What led Einstein to predict stimulated emission of photons? Discuss
briefly.
SA8.10 Which laser source uses a ceramic plasma tube and why? Give energy level
scheme of the laser ion.
SA8.11 What is the purpose of mode locking and Q-switching in laser sources?
What is the main difference in the two?
SA8.12 What is a hologram and in what respect it differs from a photograph? Briefly
explain how a hologram is made.

Note: Some of the multiple choice questions may have two or more correct
alternatives. All correct alternatives must be marked for complete answer of the
question.
MC8.1 Frequency separation between two modes of a parallel mirror resonant
cavity of length L is;
cμ c/μ cμ c/μ
(a) 2L
(b) 2L
(c) L
(d) L
ANS: (b)
MC8.2 The ratio of Einstein’s coefficients A21 /B21 is given by
8π ν 3 8π hν 3 hν 3 8π hν 3
(a) hc
(b) c
(c) 8π c
(d) c3
ANS: (d)
MC8.3 A laser beam of photons of energy 1 J hits a mirror for 30 s. Assuming
that the mirror reflects back all the photons of the beam, the momentum
transferred by the beam to the mirror is
(a) 0.1 × 10−7 kg m (b) 1 × 10−7 kg m (c) 2 × 10−7 kg m (d) 20 ×
10−3 g cm
ANS: (c) and (d)
MC8.4 A laser beam of transfer mode TM00 has beam radius of 0.01 cm at
1 m from the source and of 0.2 cm at 2 m distance from the source, the
divergence of the laser beam is;
(a) 0.5 × 10−3 rad (b) 0.5 × 10−3 degree (c) 1.0 × 10−3 rad (d) 1.0 ×
10−3 degree
ANS: (c)
MC8.5 A laser beam of wavelength 600 nm has a wavelength spread of 0.5 nm.
The frequency of the laser ν and frequency spread Δυ are respectively;
(a) 5 × 1014 Hz; 41 × 1010 Hz (b) 5 × 1010 Hz; 41 × 1014 Hz (c) 41 ×
1014 Hz; 5 × 1010 Hz (d) 50 × 1014 Hz; 4.1 × 1010 Hz
ANS: (a)
MC8.6 A laser light of frequency 5 × 1014 Hz and frequency spread 42 × 1010 Hz
is confined in a resonant cavity of 5 cm length. The number of axial modes
in the cavity is;
(a) 40 (b) 140 (c) 200 (d) 240
ANS: (b)
MC8.7 An atomic system has two excited states at energies E 1 and E 2 (E 2 > E 1 ).
If RT 1 and RT 2 denote the ratio of the number of atoms in state E 1 and E 2
at temperatures T 1 K and T 2 K respectively, then the ration R = RT 1 /RT 2
for T 2 = 2T 1 is;
( E2 −E1 ) ( E2 −E1 ) 2( E 2 −E 1 ) ( E1 −E2 )
(a) e 2kT 1 (b) e kT 1 (c) e kT 1 (d) e 2kT 1
ANS: (a)
MC8.8 Tick the correct alternative(s);
The number of stimulated photons is proportional to;
(a) Number of atoms in upper lasing level (b) number of atoms in lower
lasing level (c) number of triggering photons (d) number of atoms in the
ground state
ANS: (a) and (c)
MC8.9 In a closed tube argon ion laser the purpose of putting helium gas in the
plasma tube is to;
(a) Transfer excitation energy to argon ions (b) dissipate heat from the
central region of the plasma tube to its walls (c) facilitate plasma formation
(d) facilitate de-excitation of atoms in lower vibrational excited states of
argon ion
ANS: (a) and (d)
MC8.10 Which of the following laser source does not require any sort of pumping
(a) Nd.YAG (b) Ar ion (c) HBr (d) GaAs diode
ANS: (c)
LA8.1 What is meant by population inversion and stimulated emission? Discuss the
working of CO2 laser source and mention some of its important applications.
LA8.2 Discuss the construction and working of a semiconductor laser source. Why
it operates as a laser source only when forward current is more than the
threshold current? What is a heterojunction diode laser and why it is more
efficient?
LA8.3 What are the special characteristics of laser light? Discuss three of them in
details. Define angle of divergence of a laser source.
LA8.4 Describe Einstein’s theory of stimulation emission and calculate the ratio
of spontaneous to stimulated emission for a system in thermal equilibrium.
What are Einstein’s coefficients?
LA8.5 With the help of a neat diagram explain the working of a Nd.YAG laser. What
is the purpose of resonant cavity? Derive an expression of the threshold gain
for a parallel mirror resonator cavity.
LA8.6 Write an essay on lasers and their applications.
Chapter 9
Nanomaterials
Objective
Properties, behaviour and production of nanomaterials, along with their important
applications, will be discussed in this chapter. It is expected that a reader after going
through this chapter will be able to understand why nanomaterials are assuming added
importance and how nanomaterials with desired characteristics may be synthesised
in laboratory.
9.1 Introduction
Nanomaterials are magical materials of immense potential, though their potential has
been realised only recently. Nanomaterials belong to the world of the small and the
smallest parts, of micro- and nanotechnologies. The field of nanomaterials extends
to all branches of science; electronics, mechanics, optics, chemistry and biology, etc.
Nanomaterials are the product of nanotechnology which is the science, engineering
and technology at the nanoscale. Nanoscale ranges from 1 to 100 nanometres (1
× 10−9 –100 × 10−9 m) as shown in Fig. 9.1. One nanometre is one-billionth of
metre, a very small distance, just to have a feel of nanoscale, it may be said that our
finger nails grow about one nanometre per second, a sheet of paper is about 100,000
nanometres thick, and there are 25,400,000 nm in an inch.
The nanoworld may be considered to be intermediary between the atom and the
bulk solid. The concept of nanoworld emerged from the convergence of a mix of
scientific and technological domains which once were separate. Moreover, concept
of nano is becoming fashionable since it combines what is already known with new
concepts and gives the idea of modern technologies.
In principle nanoworld may be reached through two seemingly opposite routes,
namely the Top-down and the Bottom-up approaches. In top-down approach the
starting point is the aggregate matter (solids, liquids, gases and plasma), i.e. the
traditional large-scale world with macroscopic properties which is then analysed
https://doi.org/10.1007/978-3-031-32084-2_9
436 9 Nanomaterials
Fig. 9.1 Scales for micro- and nanoworlds
and technologically manipulated to create microsystems; like semiconductor silicon,

transistors, microprocessors, etc. The top-down route is the route of miniaturisation.
The next step in top-down process is to produce devices based on large molecules
and aggregates, i.e. nanomaterials. In Bottom-up approach, starting from billions
of atoms/molecules one selects only few molecules that follow laws of quantum
mechanics and create self-assembly of nanosystems like carbon nanotubes, etc. The
bottom-up approach may be linked to the famous lecture known as ‘Room at the
Bottom’ delivered by Prof. Richard Feynman at American Physical Society in 1959.
In this lecture Prof. Feynman raised four questions (Fig. 9.2):
(i) Why can’t materials be manipulated atom by atom?
Fig. 9.2 Top-down and

bottom-up approaches to
nanomaterials
9.2 Special Features of Nanomaterials 437
(ii) Why can’t the synthesis of individual molecule be controlled?

(iii) Why can’t all human knowledge be written on the head of a pin?
(iv) Why can’t machines to accomplish these goals be made?
Non-availability of appropriate tools and technology was the single answer to all
these questions at that time. However, at present new tools for atomic-scale char-
acterisation, new capabilities for single atom/molecule manipulation, computational
access to large systems of atoms, and convergence of almost all scientific disciplines
at nanoscale, etc. has made it possible to reach nanoworld via bottom-up route.
Various types of nanostructures which possess at least one dimension in the nanor-
ange are included in nanomaterials. Some typical examples are: Quantum dots, zero-
dimensional monostructures such as metallic, semiconductor and ceramic nanopar-
ticles of diameter 1–10 nm; one-dimensional nanostructures nano wires (diameter
of 1–100 nm), nano tubes (diameter 1–100 nm) and nano rods (diameter 1–100
nm); two-dimensional nanostructures such as thin nano films (area of several nm2 to
µm2 and thickness 1–1000 nm). Besides these individual nanostructures, ensembles
of these make high dimension arrays, assemblies and supper lattices.
SAQ: What qualifies a structure to be a nanostructure?
SAQ: What is the minimum number of atoms of average size that will make a
nanostructure?
At the nanoscale, the physical, chemical, mechanical and biological properties of
materials significantly differ in fundamental and valuable ways from the properties of
individual atoms and molecules or the bulk matte. Nanotechnology may be defined
as the research and development at the atomic, molecular and supramolecular levels,
in the length scale of approximately 1–100 nm range, through the control and manip-
ulation of matter at molecular level to design, fabricate and use materials, devices
and systems with fundamentally new properties and functions because of their small
structure.
9.2 Special Features of Nanomaterials
Nanomaterials exhibit some special properties that are not found in normal macro or
even micro-materials. Most of these special features of nanomaterials arise because
of (a) large surface to volume ratio, (b) large fraction of surface atoms, (c) high
surface energy, (d) spatial confinement and (e) reduced imperfections. Some of
these properties are discussed below.
(i) Large Surface Area
Surface-to-volume ratio, i.e. the surface area for a given volume is much larger in case
of nanomaterials as compared to the similar volume of ordinary largescale or even
microscale materials. Importance of available surface area may be easily understood
by the example of chocolate. If one puts a block of chocolate of size 3.2 × 2.3 ×
438 9 Nanomaterials
1.3 cm (of surface area approximately 26 cm2 ) in his mouth; only 26 cm2 area will
come in contact with taste buds of the mouth. However if the block of chocolate is
broken into two equal pieces and the two pieces are put in mouth roughly 31 cm2
area will be available for the taste buds to enjoy chocolate. Division of chocolate
block in smaller pieces increases the available area so much so that if it is broken
into 1 nm cubes, the available surface area will be around 500,000,000 cm2 , about 10
football fields. In electric batteries current producing chemical reaction takes place
at the surface of electrodes. Therefore, electrodes made of nanomaterials expose
more surface area and thus have high current capacity. Large number of chemical
and physical processes takes place at the exposed surface, for example, adsorption of
gases, catalyst actions, etc. proceed at the exposed surface. Nanomaterials, therefore,
play important role in such processes, metallic nanomaterials are used as very active
catalysts. Chemical sensors made from nanoparticles and nano wires enhance the
sensitivity and sensor selectivity.
SAQ: Surface-to-mass ratio of nanostructures will be large or small?
(ii) Colour of Light-Emitted/Absorbed Depends on the Size of the Nanostruc-

ture
Optical properties of nanomaterials are very fascinating and useful. Optical properties
of nanomaterials depend on nanostructure size, surface characteristics, shape, and
other parameters including the level of doping and interaction with surrounding mate-
rials. Many devices based on special optical properties of nanomaterials have been
made, for example, optical detectors, lasers, sensors, imaging devices, phosphor, and
photocells, along with devices used in photoelectro-chemistry and in biomedicine.
When ultraviolet light is made to fall on a nanostructure, the light emitted (fluo-
rescence emission) by nanostructure depends on the size of the structure; Quantum
dots 3.2 nm in diameter emit blue light, while 5 nm diameter quantum dots emit red
light. Fluorescence emission of light of different colours ranging from blue to red
occurs by quantum dots of sizes ranging from 3.2 to 5.5 nm as shown in Fig. 9.3.
This is due to the fact that with the increase of quantum dot size the number of atoms
in it also increases which in turn reduces the average frequency or the average energy
of emitted photons.
Similarly, the characteristic absorption spectra frequency of gold nanoparticles
and nanorods depend on their sizes as shown in Fig. 9.4. Sketches (a) and (b) in
Fig. 9.4 corresponds to the gold quantum dots of size 15 nm and 30 nm respectively,
it may be observed in this figure that the characteristic absorption frequency does not
change much with the size of the quantum dot. However, the change in characteristic
absorption frequency becomes quite appreciable when quantum dot is replaced by a
nanorod of length 2.5 nm (curve c) and of length 7.5 nm (curve d).
White light falling on nanostructures may interact in different ways with the
nanomaterials producing colour effects. For example if the incident light undergoes
interference on falling on nanomaterial, the colours corresponding to constructive
interference of waves will appear in the output. This happens when white light
falls on the wings of a butterfly. Butterfly wings have photonic crystals as natural
Fig. 9.3 Light emitted by quantum dots depends on their size
Fig. 9.4 Typical absorption spectra of quantum dots and nano tubes of different dimensions
nanomaterials and interference of light scattered from different nanophotonic crystals

produce different colours.
In case the incident white light is simply scattered by the nanostructures, the
scattered light may have several colours, since the colours of the scattered light
depends on the size of the scatterer.
When light hits a metal surface of any size some of the light waves propagate along
the metal dielectric interface giving rise to the phenomenon of Surface Plasmon,
where some of the conduction electrons of the metal move parallel to the interface.
In case of bulk metals the plasmon electrons move freely and do not produce any
special effect. However, in nanoparticles the surface plasmon is localised in space
and oscillates back and forth in a synchronised way in a small space. This is called
440 9 Nanomaterials
Localised Surface Plasmon Resonance (LSPR). When the oscillation frequency

of the surface Plasmon matches the frequency of the incident light the process gives
rise to resonance, resulting in strong absorption of incident light. As a result quantum
dots of metals like Gold and Silver because of LSPR may display colours which do
not appear on scattering with their bulk materials.
(iii) Enhancement in Mechanical Properties of Nanomaterials
Many of the mechanical properties like, hardness, fracture toughness, elastic
modulus, fatigue strength and scratch resistance, etc. of nanomaterials are found
to be enhanced as compared to the corresponding bulk materials. This happens
primarily because of the two reasons; (i) structural perfection (very few deformities/
dislocations/defects/micro twins/impurity precipitate, etc. are present in nanostruc-
tures) of nanostructures and (ii) the self-purification property of nanomaterials.
Imperfections/impurities in nanomaterials are highly energetic, which migrate to
the surface to relax themselves under annealing/purification of the material. As such,
both the interior and the outer surfaces of nanomaterials are highly defects free
which do not yield at low pressures. Pure nearly spherical nanosilicon particles of
diameters from 20 to 50 nm may withstand pressures up to 50 GPa which is almost
4 times greater than for bulk silicon.
High hardness is one of the novel properties of nanomaterials. A variety of super
hard nanocomposites have been made from nitrides, borides and carbides using
plasma-induced chemical and physical vapour deposition. The hardness of some
of Si3 and N4 nanocomposites exceed 50 GPa which is almost double of 20 GPa,
the hardness of corresponding normal composites. Nanocomposites with enhanced
surface hardness may be used for reinforcement of polymer materials as light weight
and high-strength materials, flexible conductive coatings, wear resistant coatings and
as tougher and harder cutting tools, etc.
Carbon nanotubes discovered by electron microscopist Sumio Iijima of NEC
laboratory of Japan in 1991 gave a big boost to research on nanostructures. Carbon
nanotubes/fibers are found to have novel mechanical properties. On account of their
innumerable applications we shall talk about carbon nanotubes in more details later;
however, it may be mentioned that the measured value of Young’s modulus for single-
walled carbon nanotubes lies between 0.5 and 5.0 TPa, which is much higher than
the value for high-strength steel which is around 0.2 TPa.
SAQ: LSPR effect does not occur in bulk material or it is not observable in bulk
material?
SAQ: Why nanostructures have almost no defects?
(iv) Electronic Properties of Nanomaterials
Bulk materials, on the basis of their electrical properties, may be categorised as

conductor, semiconductor and insulator. In case of solids, the magnitude of forbidden
energy gap or energy band gap (E g ) decides the electronic/electrical nature of the
material; E g ≈ 0 defines a conductor; E g ≈ few eV (up to 5 eV) are semiconductors
and E g > 7 eV (say) a are insulators. Band gap energy for semiconductor increases in
Fig. 9.5 Energy states of a

single atom and energy band
gaps for bulk and
nanostructures of a
semiconductor
nanostructure as compared to the bulk material because of the quantum confinement

in case of nanostructures. This happens because of the loss of energy states due to
small size and much less number of atoms as shown in Fig. 9.5 where E g > E g2 >
E g1 , where E g , E g1 and E g2 are respectively the energy band gaps for an atom, bulk
material and nanomaterials.
In case of bulk semiconductors which at T > 0 K may contain large number
of delocalised electrons in conduction band and holes in valence band, there is no
confinement of electrons/holes in any dimension. Both the conduction and the valence
bands, in a bulk semiconductor, have large number of closely spaced energy states
for these particles. These energy states are so closed to each other in energy that it
becomes impossible to talk about a specific energy state at some energy E. Instead
one talks about the level density, the number of levels per unit volume within energy
E and (E + dE). The state density of allowed energy states in conduction band may
be denoted by N(E) and holes in valence band by P(E). In order to calculate the
concentration (number density) of electrons in the conduction band of a semicon-
ductor we chose a small energy interval dE around energy E in the conduction band
and denote the concentration of electrons in this energy interval by ne (E). Now the
magnitude of ne (E) is given by the product of two factors: (i) the density of allowed
energy states at energy E, N(E) and (ii) the probability F(E), that this energy range
is occupied by electrons. Thus
n e (E) = N (E)F(E) (9.1)
To obtain the total number of conduction electron ne in the complete conduction

band, one has to integrate the above expression from the energy E bC at the bottom
of the conduction band up to the top E top , i.e.
Etop Etop
ne = n e (E)dE = N (E)F(E)dE (9.2)
E bc E bc
Theoretical value for N(E), density of state at energy E and (E + dE), (assuming
free electron to be a particle in box), is given by
442 9 Nanomaterials
3/2
2m e
N (E) = 4π E 1/2 (9.3)
h2
And F(E), the probability that the state at energy E is occupied by electron is
given by quantum statistics as;
1
F(E) = ( E−Ef )
(9.4)
1+e kβ T
Concentration ne of electrons in conduction band and np of holes in valence band

govern the current flow and electronic behaviour of buck semiconductor. The varia-
tion of N(E) with E, the density of states in conduction band with energy E, has the
shape shown in Fig. 9.6a for bulk semiconductor. However, the energy dependence
of state density for quantum dot (confinement in all three dimensions), quantum wire
(confinement in two dimensions) and quantum well (confinement in one dimension)
are shown in figures (b), (c) and (d) of Fig. 9.6. Since electronic properties of materials
depend on their charge carrier state densities, the electronic behaviour of nanomate-
rials is very different from bulk material as well as from each other and depends on
the dimensions/shape of the nanostructure.
It is obvious that only those properties of the matter that depend on the number of
atoms in the given specimen will drastically change in going from bulk to nanostruc-
tures. For example in conductors, like metals, the Fermi level lies in the overlapping
region of conduction and valence bands while the conduction band is generally half
filled (at room temperature). Further, the density of unoccupied levels in conduction
band is quite large. This characteristic, large number of unoccupied levels does not
depend too much on the number of atoms in the specimen, till the number of atoms
in the specimen goes down below 100 atoms. If the size of nanometal particle is
Fig. 9.6 Electron state

densities for, a bulk
semiconductor, b quantum
dot, c quantum wire and
d quantum well
reduced sufficiently, then the continuous density of electronic states is broken into
discrete energy levels. The spacing ΔE between energy levels depends on the Fermi
energy E f of the metal and the number of electrons N in the metal specimen as given
by
4E f
ΔE = (9.5)
3N
Fermi energy E f for most metals is of the order of 5 eV. Discrete electronic levels
in gold nanoparticles have been observed experimentally in far-infrared absorption
measurements confirming expression (9.5).
Solids may be classified as conductors, semiconductors and insulators according
to the value of their resistivity, ρ, (resistance between the opposite faces of a unit cube
of the material). Resistivity arises essentially because of the scattering of delocalised
electrons by crystal lattice while moving under the applied electric field. According
to Drude theory, at room temperature there is some number of delocalised electrons in
conductors which are moving with thermal speeds (∼ 105 m/s) in random directions.
When a voltage V is applied across a conductor of length L, an electric field ∈ (= V /
L) is established in the specimen which forces the delocalised electrons in it to move
opposite to the direction of the field ∈. The electric field ∈ exerts a force on each
free electron which imparts an additional component of velocity called drift velocity
velocity V d to each free electron. Electrons moving under the influence of the electric
field undergo frequent collisions with the crystal lattice which randomise the motion
of drifting electrons. As such scattering by crystal lattice may be treated like a friction
force which counterbalances the force due to the electric field. This results in Ohm’s
law V = RI. Scattering events have a mean free path λ, the average distance travelled
by an electron between two successive scattering events, which is of the order of nm
in many materials at room temperature. If the size of a nanostructure is of the same
order of magnitude as the mean free path λ, then Ohm’s law may not be obeyed. In
such a situation the electron transport process is totally quantised and current–voltage
relationship may be quite different than Ohm’s law.
(v) Thermal behaviour of nanostructures
Many properties of nanostructures like, mechanical, electronic and optical behaviour
have been well studied in recent past. However, thermal behaviour of matter at
nanoscale could not be investigated in details till now. There are several prob-
lems in this regard. Firstly, the definition of temperature at nanoscale is itself quite
ambiguous. In non-metallic materials heat energy is generally transported through
photons and or phonons, both of which may have widely different frequencies and
mean-free-paths (mfp). Heat-carrying photons often have wavelengths and mfp in the
range of nanometres at room temperature. Thus nanostructure size and heat photon
parameters are of the same order of magnitudes. A temperature is defined only when
the system is in equilibrium. In bulk material it is possible to define local tempera-
tures in different small regions which may be treated as if they are in steady state/
equilibrium, and thus it becomes possible to study the process of heat transport based
444 9 Nanomaterials
on the temperature distribution in the material. In case of nanostructures, it is not

possible to subdivide the structure in still smaller regions and achieve thermal steady
states in these sub-nanoregions. Theoretical study of heat transport in nanostructures
is equally challenging. The three fundamental theoretical approaches of thermal
transport, namely numerical solutions of Fourier’s law, molecular dynamic simu-
lations (MDS) and calculations based on Boltzmann transport equations, all have
limitations in case of nanostructures. In spite of all these problems, recent advance-
ments and introduction of atomic force microscopy (AFM) have shown that some
nanomaterials have extraordinary thermal properties as compared to their macro-
scopic counterparts. For example, it has been found that silicon nano wires have
much smaller value for thermal conductivity as compared to the bulk silicon; and
carbon nanotubes have a much larger value of thermal conductivity along the axial
direction as compared to the other directions. Therefore nanotubes, in general, show
high anisotropy in heat transport. Interfaces within the nanostructures are found to
impede the transport of heat, perhaps because of the scattering of heat carrier photons
at interface boundary.
Carbon nanotubes made up of seamless graphitic cylinders are expected to
have extraordinarily high thermal conductivities, rigidity and virtual absence of
any defects makes these tubes the front running candidate for efficient thermal
conductors. Theoretical studies on temperature dependence of thermal conductivity
of carbon nanotubes show that thermal
W conductivity has a peak value of around
4 × 104 Watt per metre per Kelvin m·K at around 100 K and then decreases almost
exponentially, as shown in Fig. 9.7.
Nanofluids are also being looked as promising candidates for future thermal trans-
port systems. Nanofluids are solid–liquid composites made by mixing/suspending
in small percentage by volume nanoparticles of size in the range 1–100 nm in an
appropriate liquid. Many different kinds of nanoparticles like nanoparticles of metals,
oxides, carbides, nitrites, and nanofibers may be used. Liquids like water, oils, ethy-
lene glycol, etc. may be used as the base fluid in which above-mentioned nanostruc-
tures are dispersed to make nanofluid of desired properties. The thermal conductivity
of nanofluids is found to have much higher value as compared to the base fluid further;
the change in in thermal conductivity is found to be a function of the percentage of
Fig. 9.7 Temperature

dependence of thermal
conductivity of carbon
nanotubes based on
theoretical calculations
nanomaterials by volume. For example in some experiments where nanoparticles

of Al2 O3 were mixed (≈ 4.5%by volume) with water, the thermal conductivity of
the nanofluid got enhanced by about 30%. Nanofluid made by suspending carbon
nanotubes in oils, 1% by volume, is found to have thermal conductivity almost 160%
of the base oil.
(vi) Magnetic behaviour of nanomaterials
Magnetic nanoparticles have been found in some rocks and have been used to deter-
mine the Earth’s magnetic field, i.e. its strength and direction. Some birds and living
creatures are found to have clusters of nanoparticles of 2–4 nm size in specific areas
of their bodies, for example in their beak area, which is believed to help them in
navigating and with their homing ability. It is interesting to note that even in bulk
magnetic materials some structures may be in the nanoscale; for example domain
walls in a ferromagnetic material have the dimensions of the order of 60 nm. More-
over, domains that are near to the grain boundaries or the surface of the macrostructure
are themselves nanostructures.
Magnetic behaviour of materials is generally studied in terms of the magnetic
moment per unit volume or magnetic moment per atom. Experiments have shown
that magnetic moment per atom increases as the dimensionality of the system is
decreased which is indicated by the following data (Table 9.1) on iron and nickel
structures.
It is interesting to note that the experimental data on magnetic moment per atom
indicated above goes against the expected decrease of it with the reduction in dimen-
sionality. It is because the surface spins are generally not ordered in the same direction
(as compared to the spins in the interior) and with the decrease of dimensionality the
surface area increase; hence, magnetic moment per atom should decrease with the
reduction of dimensionality.
In bulk ferromagnetic material there are several domains that contain aligned
magnetic dipoles. With the reduction of size of the ferromagnetic specimen the
number of domains in it also decreases. The magnetic energy of a ferromagnetic
particle of radius r varies as r 3 , while the domain wall energy goes as r 2 . Therefore,
there is a critical value of the particle radius rc below which domain walls are not
sustained by the particle. As a result nanosize particles may have alignment of all
magnetic dipoles in one direction giving rise to supermagnetism. Collapsing of
domain walls with reduction in the dimension of the nanostructure may convert
a nanoparticle of antiferromagnetic and ferrimagnetic materials into ferromagnetic
material.
Table 9.1 Dimensionality dependence of magnetic moment per atom

Material Magnetic moment per atom in unit of Bohr magneton per atom (µB /atom)
Bulk Two dimensional (2 D) One dimensional (1 D) Quantum dot (0D)
Iron 2.27 2.96 3.30 4.00
Nickel 0.56 0.68 1.10 2.00
446 9 Nanomaterials
Measurement of magnetic properties of nanostructures is a challenging task as

magnetic properties depend strongly on the number of atoms in the nanostructure.
Coagulation or contamination of nanostructure either during production or during
measurement may alter the magnetic behaviour. Temperature also plays a crucial role
and controlling it during measurement poses big problem. Since most of the atoms
in nanoscale structures reside at the surface, surface effects dominate the magnetic
behaviour and therefore, due care be taking in developing an appropriate theory to
compare the measured magnetic properties of nanostructures.
SAQ: What is supermagnetism and why is it likely to be observed in nanostructures?
9.3 Technology Used for the Study of Nanostructures
Tools/microscopes based on three different technologies are generally used to study

the structural details of nanomaterials.
(a) Scanning Tunnelling Microscope (STM) Scanning tunnelling microscope was
invented by Gerd Binning and Heinrich Rohrer in 1981, for which they got Nobel
Prize in 1986. The instrument works by scanning a very sharp metal wire tip
over the surface of the specimen to be studied. The wire tip is kept very close to
the surface and an electric potential difference is maintained between the wire
tip and the surface. The applied potential difference between the wire tip and
the specimen surface produce field emission of electrons which tunnel through
the gap between the tip and surface constituting a tunnelling current (Fig. 9.8).
Characteristics of tunnelling current depend on the surface structure of the
specimen, and thus, it works as an ‘eye’ to decipher the surface topology/
structure of specimen. Precise angstrom level control on tip position is achieved
using piezoelectric effect controlled motors. An electronic feedback loop in
the tunnelling current circuit produces a 3-D image of the scanned structure.
Tunnelling current essentially depends on the gap between the metallic tip and
specimen surface and fluctuates with the motion of the wire tip as the gap
changes on account of irregularities in specimen surface. The feedback loop
Fig. 9.8 Lay out of scanning

tunnelling microscope
9.3 Technology Used for the Study of Nanostructures 447
shifts the metallic tip in desired direction so as to maintain the same gap and
tunnelling current. The amount of shifts in the position of the tip to keep constant
tunnelling current (produced by the feedback loop) are recorded in a computer
and are used to produce a 3-D surface image by the computer software. This
method of using (STM) is called constant current method. In case the specimen
surface is rather smooth, the feedback loop may be removed and the tunnelling
current reading itself may be converted into a 3-D surface image or surface
topology.
The biggest drawback of STM is that it can record the surface topology only
of those surfaces which are good conductor of electricity. Tunnelling current is
quite small and in STM it flows through the surface under study. In case the
surface offers high resistance it may not be possible to record small changes in
tunnelling current accurately.
(b) Atomic Force Microscope (AFM) The STM may decipher the surface topology
only of those surfaces that are good conductor of electricity. Atomic force
microscopy was developed to overcome this drawback of STM. AFM has the
advantage that it may image the topography of any type of surface, including
ceramic, polymer, composite, glass and of biological samples. AFM was
invented by Binning, Quate and Gerber in 1985.
In AFM a very sharp tip is attached at one end of a metallic strip, the other
end of which is clamped in a holder. Metallic strip, clamped at one end and
holding a fine tip at the other free end, works like a cantilever. The strip holder
may be moved in X-Y plane with the help of synchronised motors, and the X, Y
coordinates of the tip may be recorded with precision. The pointed tip is kept
very close to the surface under study (like that of STM); however, no current
is made to pass between the tip and the surface in case of AFM. In this case
atomic force between the tip and the surface either pull the tip towards surface
or push the tip away from the surface. Thus, the tip which is suspended by a
metallic cantilever, moves in the Z-direction, towards the surface or away from
it when (the tip is) scanned over the surface. Atomic force between the tip and
the surface depends on the distance of separation between them. If there is a
depression in the surface at some point, the tip will experience a smaller atomic
force and will move up while at a point of bulge the separation between tip
and surface will become less, atomic force will increase and the tip will move
downwards. Thus displacement of the sharp tip in Z-direction may be converted
into the topology of the surface.
Schematic layout of atomic force microscope is shown in Fig. 9.9. Tip
displacement in Z-direction is recorded using a laser beam which is made to
hit the top of the sharp tip. The laser beam reflected from the top of the sharp
tip is recorded in a position sensitive laser detector. In their original instrument,
Binning et al. used diamond for sharp tip and gold foil as cantilever strip.
Atomic force microscope relies on force between the tip and the sample
surface. Atomic force and distance (of the tip from the surface) curve is shown
in Fig. 9.10. Atomic force is not measured directly but is calculated by measuring
the deflection of the tip knowing the stiffness of the cantilever strip using the
448 9 Nanomaterials
Fig. 9.9 Schematic diagram of atomic force microscope
Fig. 9.10 Force–distance

Repulsive
curve for AFM Tip hard pressed:

repulsive force Tip far from the surface
no deﬂection
Atomic force
C A
Attractive
Tip pulled towards surface:

B attractive force
Range of
atomic
force
0.10 nm 100 nm
Tip distance from the surface
Hook’s law; F = − kZ . Here k is the stiffness constant of the strip and Z the
displacement of the tip.
The range of atomic force extends from about 0.10 to 100 nm. A typical curve
showing the variation of atomic force with the tip-surface distance is shown in
Fig. 9.10. When the tip of the AFM is more than 100 nm away from the surface
(on the right side of point A in the figure) it does not experience any atomic force.
On moving the tip towards left of point A, the tip experiences an attractive atomic
force that pulls the tip down and the maximum attractive force is experienced
at the distance corresponding to point B. The exact distance corresponding
to point B is different for different surfaces, depending on the nature of the
surface, i.e. the atoms/molecules of the surfaces. On further reducing the relative
separation between the tip and the surface beyond point B, the atomic attractive
force diminishes, reducing the downward displacement of the tip. Ultimately,
at separation corresponding to point C atomic force experienced by the tip
becomes zero again and the tip assumes its undeflected position. When tip-
surface distance is reduced beyond point C, which means that the tip is hard
pressed into the surface it experiences a force of repulsion and the tip is pushed
back. In this way, sensing the deflection of the tip, the detailed structure/topology
of any surface, conducting or non-conducting may be studied using an atomic

force microscope.
(c) Transmission Electron Microscope (TEM) In optics we have studied that the
spatial resolution of an image depends on the wavelength of the light that formed
the image, shorter the wavelength better is the resolution. It is also known that
a moving particle has an associated wave the wavelength of which depends
inversely on the mass and the speed of motion of the particle. Electrons of mass
9.1 × 10−31 kg moving with speeds as high as 0.5c (c is the speed of light) carry
a wave of very small wavelength. Images formed by the scattering of electron
waves have very high resolution. This principle is used in transmission electron
microscopy. TEM is used to study the fine structural details of the inside of
micro and/or nanoparticles, sometimes up to the level of an atom. TEM may
have an amplification factor of 2 × 103 , i.e. the image of an object may be two
million time larger than the object.
Working of (TEM) may be understood with reference to Fig. 9.11. The elec-
tron gun at the top provides a beam of high-energy electrons. Electron beam
is collimated to a fine pencil using two electromagnetic collimation condenser
lenses. There is an aperture control system after each lens. Collimated beam of
high energy electrons falls on the specimen material held on specimen holder.
Electron beam after scattering and or partial absorption by internal features of
the specimen passes through the electromagnetic objective lens and produces
the amplified image of internal structures of the specimen on a fluorescence
screen for visual observation. The image may also be photographed with the
help of a camera.
Fig. 9.11 Schematic

diagram of transmission
electron microscope
450 9 Nanomaterials
(d) Optical Tweezers (OT) Optical tweezers (OT) are instruments based on intense
laser beams that may hold or trap particle from the size of an atom up to strand of
DNA and living cell. In most cases the desired particle is trapped at a particular
point in a plane called the focal plane. Optical tweezers was invented by Arthur
Ashkin for trapping essentially biological molecules, cell, etc. for which he
received the Nobel Prize in 2018.
Using only light (laser beam) OT is able to hold or influence the motion
of objects (nano to microsize including biological cells, etc.) in a non-contact
way. This makes OT especially useful for holding and studying nanostructures
which are difficult to manipulate using conventional means such as mechanical
tweezers or micropipettes.
Most optical tweezers use a highly focused and intense laser beam, usually
of the wavelength in the range of 0.5–1.0 µm (visible to near infrared). The
laser beam is often focused using a microscope, the objective lens of which
focuses the beam at a small spot in the focal plane where the sample having
nanostructure to be studied is kept. The OT holds the nano or the other desired
microstructure (present in the sample) at the focal spot of the laser beam. The
laser beams used in OT have a Gaussian intensity profile, i.e. the intensity of
the beam is very high in the central region and decreases towards the periphery.
In order to understand the operation of the OT, let us consider a very small
(nano or microsized) particle of some dielectric medium which is lying at some
location in the laser beam. This dielectric particle is often called a bead. The
bead in the path of the laser beam may experience three different forces due
to three interactions it may have with the beam. These three interactions may
be understood in terms of the photon nature of laser beam. The laser beam
of frequency ν may be considered as a collection of photons, each of energy
E = hν. Now a photon of energy E carries a linear momentum = Ec = hν C
, where
c is the velocity of light. The laser beam photons may be scattered (reflected) by
the surface of the bead giving rise to what is called the Scattering force. The
scattering force pushes the bead towards the focal spot in the focal plane. Another
force, called the Gradient force, comes into play because of the intensity profile
of the laser beam. Since the central part of the laser beam carry larger number
of photons per unit volume as compared to the number density of photons at
periphery, a net gradient force is applied to the bead which pushes the bead
towards the centre of the beam. When scattering and gradient forces are larger
than other forces acting on the bead, like the force of gravity and the force of
Brownian motion, the bead is steered towards the focal spot and is held there. A
third type of force called Absorption force may come into play because of the
absorption of the laser radiations by the bead. Absorption force typically behaves
like the scattering force; however, too much absorption of laser radiations may
increase the temperature of the bead (or the specimen) which may be harmful
in case of biological samples.
A typical layout of an optical tweezers is shown in Fig. 9.12 where a laser
source produces an intense beam which is first expanded and then steered using
a mirror and a halfwave plate (to compensate for phase change) to a microscope
Fig. 9.12 A typical layout

of an optical tweezers
which tightly focuses the beam to the focal spot. The sample containing nano/
microstructure is placed in the focal plane of the microscope objective. An
eyepiece and camera are provided with the set-up to see and record the event.
In order to understand the origin of scattering force it may be mentioned that
each incident laser beam photon scattered from the surface of the bead imparts
a linear momentum or apply a force (called the recoil force) to the bead. For an
isotropic scatter, the resulting forces cancel in all directions except the incident
direction. Thus, as a result of scattering (reflection) of incident laser photons
from all sides of the bead surface in all possible directions, a net force pushing
the bead in the direction of beam propagation is generated which is termed as
the Scattering force. Figure 9.13 shows a bead placed in a laser beam, the
surface of the bead is bombarded with laser beam photons from all directions.
Let us pick a typical photon which is scattered at point A of the surface. Since
after scattering the momentum of the scattered photon has changed, it exerts a
force say, F 1 on the bead to recoil back. Now force F 1 may be resolved into two
perpendicular components F 1 H and F 1 V . Similarly, other photons on scattering
will also apply forces F 2 , F 3 ,…F N etc., on the bead each of which may be
decomposed into two components F 2 H , F 2 V ; F 3 H , F 3 V …, F N H , F N V … etc. It
can be shown that in case∑ the scattering is isotropic, the sum of F N H components
will add up to zero, i.e. FNH = 0. The y-components of all forces add up to
give a non-zero
∑ V value of force pointing in the incident direction; i.e. Scattering
Force = FN.
Gradient force is the result of Gaussian intensity profile of the laser beam.
The intensity of cylindrical laser beam is a maximum at the axis and decreases
as one move away from the centre. Now laser beam being an electromagnetic
wave carry both electric and magnetic fields. The intensity of the electric field
produced by laser beam photons is a maximum at the beam axis and decreases
towards the periphery. Further, the electric field associated with photons induces
fluctuating electric dipole in the dielectric material of the bead. The fluctuating
dipole in the bead interacts with inhomogeneous electric field of incident beam
and experience a force towards the stronger region of the electric field. Since
the force is proportional to the electric field gradient, it is a maximum at the
452 9 Nanomaterials
Fig. 9.13 Origin of

scattering force
outer rim of the beam and is a minimum at the central region, at the beam axis.
Thus, the gradient field works as a spring which always tries to keep dielectric
bead at the centre of the focal spot. It is interesting to note that the magnitudes
of both Scattering force and the Gradient force are of the order of 10−12 N
(pico-Newton) which are enough to hold the nano or microstructures fixed at
the focal point.
In case the bead is transparent to the laser beam, the action of the gradient
force may be explained through Fig. 9.14. Let the bead be away from the focal
spot and placed such that a part of it lies in the stronger and other part in the
weaker sections of the beam intensity. The laser rays from the stronger central
part are shown with broad red arrow (as they will be large in number per unit
volume of the beam) while corresponding rays from the weaker part of the beam
by narrower purple arrow. The refracted rays are also indicated in the figure.
Both refracted rays will exert recoil forces F 1 and F 2 on the bead as indicated
in Fig. 9.14. Forces F 1 and F 2 can be resolved into X- and Y-components. The
vertical components F 1 V and F 2 V will add up to give the Scattering force that
steers the bead in the incident direction. The horizontal components F 1 H and
F 2 H will oppose each other; since force F 1 >> F 2 (because of the intensity
profile of laser beam), the net horizontal force on the bead will be (F 1 H − F 2 H )
that is the gradient force which tries to bring the bead at the axis of laser beam.
As such under the joint action of the scattering and the gradient forces the bead
is held at the focal spot.
SAQ: What will happen to the gradient force if the laser beam intensity is least at
the axis and increases with distance from the axis?
9.4 Techniques of Producing Nanostructures
In spite of it being a relatively recent field of research, several techniques for

producing/fabricating nanostructures have developed in this short time. Nanostruc-
ture fabricating techniques, as usual, may be divided into two categories; bottom-up
techniques where the desired nanostructure is synthesised by picking up individual
atom/molecule from different or a given bulk material. It is like building a house
9.4 Techniques of Producing Nanostructures 453
Fig. 9.14 Action of gradient

force
brick by brick. The other technique, classified as top-down technique, is based on

taking a specimen of bulk material and removing material from the specimen in
several steps to reach the desired nanostructure. It is like taking a mono block of a
stone and to chisel out parts of the stone to make a sculpture. It is surprising that this
top-down technique was used in ancient India to fabricate big temple complexes from
a single granite mountain block; as an example the Shiva temple of Ellora Caves.
Some important fabrication techniques both from bottom-up and top-down methods
will be outlined in the following.
9.4.1 Bottom-Up Techniques
Depending on the physical phase of the reactants, the bottom-up technique may be
further divided into two groups: (i) the gas-phase methods and (ii) the liquid-phase
methods.
(i) Gas phase Methods
(a) Chemical Vapour Deposition (CVD) Chemical vapour deposition (CVD)
is a versatile process in which gas-phase molecules are decomposed to reac-
tive species leading to the film or the particle growth. This method may be
used to deposit a wide range of conducting, semiconducting and insulating
materials. Recently the method has been used for controlled production of
nanomaterials in porous hosts. The two basic advantages of (CVD) tech-
nique are: (i) the ability to controllably create films of widely varying stoi-
chiometry that is films containing constituent materials in different ratio
and (ii) to uniformly deposit thin films of materials, even onto nonuni-
form shapes. The layout of the experimental setup generally used in (CVD)
method is shown in Fig. 9.15, where the 3D-substrate sample on which thin
454 9 Nanomaterials
film of the nanostructure is to be deposited is kept on a quartz boat in a

quartz tube. There are inlet and outlet ports at the two ends of the quartz
tube for the desired vapour to pass through it. The quartz tube is kept in an
oven the temperature of which may be controlled to any desired value.
The deposition process occurs in three successive stages: (a) introduction
of the volatile precursor, (b) adsorption of precursor vapours on the sample
surface and (c) decomposition of these products on the heated sample
followed by nucleation and growth of the solid layer/grains, followed by
the formation of volatile by-products and their removal.
(b) Plasma Arcing Plasma arcing is the method by which fullerenes and
carbon nanotubes were fabricated in sufficient amount for the first time.
The method employs two graphite electrodes kept some distance apart in a
chamber which is filled with some inert gas, mostly helium, at low pressure.
When a potential difference is established across the two electrodes, at a
certain voltage an electric arc discharge takes place, the inert gas between
electrodes ionises forming plasma in the electrode gap. As a result of plasma
arcing, cations are evaporated from the Anode and get deposited on the
cathode in the form of soot and in rod-like morphology. The soot and rod-
like deposits on cathode contain in sufficient amount fullerenes and different
carbon nanostructures including carbon nanotubes (Fig. 9.16).
Fig. 9.15 Layout of CVD

system
Fig. 9.16 Schematic

diagram of plasma arc
assembly
The production of carbon nanostructures, particularly of nanotubes,

strongly depends on the uniformity of plasma arc and the temperature of the
evaporated material deposited on cathode. In order to make uniform plasma,
the anode of the assembly is rotated using a motor. Inert gas pressure also
plays important role in determining the yield of nanostructures.
SAQ: What are the main advantages of CVD method over the plasma arc method?
(ii) Liquid Phase Methods

(c) Sol-gel Synthesis This is one of the frequently used methods of bottom-up
technique using synthesising material(s) in liquid phase. It is used for the
synthesis of various nanostructures, especially metal oxide nanoparticles,
nanostructured surfaces and nano-three-dimensional objects. In the first
step of the process a Sol of the precursor, which are usually metal alkox-
ides, is made. Sol is a colloidal solution of alkoxide and some solvent,
which may be water or alcohol. In colloidal solution particles of metal
alkoxides remain suspended in the solvent. Next step is the preparation
of a uniform gel of the molecular precursor by heating and stirring. The
colloidal suspension on heating and stirring evolves forming network in
a continuous liquid phase by the process called gelation. Gelation which
is a slow chemical process of hydrolysis or alcoysis may be enhanced by
using some catalyst. Gel formed either by the hydrolysis or alcolysis of the
alkoxide is wet or damp and has to be dried by some appropriate method
depending on the desired properties of the gel. Dried gel is then powdered
and calcined. Drying, which essentially means removal of solvent from
the gel structure, is a crucial step. The ‘removing of solvent’ method is
selected according to the application for which the gel is to be used. Gels,
dried in different ways contain metal particles in different nanostructural
forms, are directly used in industry for surface coating, building insulations
and producing special cloths etc. Sol-gel method is cost-effective, and due
to low reaction temperature there is excellent control over the chemical
composition of the products. Further, grinding the dried gel by special
mills it is possible to produce nanoparticles. Sol-gel method is capable of
producing high quality nanoparticles of same size on an industrial scale.
This method can also create two or more types of nanoparticles simultane-
ously, which means that alloy products may be synthesised in one step by
mixing two or more metal (or metal oxides) in certain proportion. Steps
of the sol-gel methods are shown in Fig. 9.17.
(iii) Solid and Liquid Phase Methods
(d) Self-Assembly Self-assembly is the spontaneous molecular arrangement
of the disordered entities of molecules into ordered structures resulting
from specific local interactions among the components themselves. Most
of the biological nanostructures, like construction of cell membranes,
helical structure of DNA, folding of polypeptide chains, etc., are the
456 9 Nanomaterials
Fig. 9.17 Steps of Sol-gel method
outcome of self-assembly. It also accounts for the construction of molec-

ular crystals, phase separated polymers, colloids, etc. The process of self-
assembly plays a key role in design, synthesis and development of new
nanomaterials.
Self-assembly, on the basis of the size and the nature of the building
blocks, may be divided into three classes: atomic, molecular and
colloidal self-assemblies. Similarly, on the basis of the systems and where
it occurs, self-assembly may be classified as biological or interfacial.
Further, depending on the process responsible for self-assembly, it may be
called thermodynamic self-assembly or kinetic self-assembly. Atomic,
molecular, biological and interfacial self-assemblies come under thermo-
dynamic class, while colloidal self-assemblies and some interfacial self-
assemblies fall under kinetic class. Self-assemblies may be directional
or random; atomic and biological self-assemblies are directional while
others like colloidal, molecular and interfacial self-assemblies are random.
The interactions that are involved in self assembly process are weak,
non-covalent, linked via van der Waals forces, hydrophobic, electrostatic
and hydrogen bonding which are normally weak individually in compar-
ison to covalent bonding but when present in large number they form very
stable self-assembled structures. Typical examples and the size scales for
atomic, molecular and colloidal self-assemblies are shown in Fig. 9.18.
9.4.2 Top-Down Techniques of Fabricating Nanostructures
In top-down technique of fabricating nanostructures, material is selectively removed

from a given precursor material to achieve the desired nanostructure. Following
methods are generally used for making nanostructures: (i) mechanical milling, (ii)
Fig. 9.18 Typical examples of atomic, molecular and colloidal self-assemblies
laser ablation, (iii) nanolithography and etching, (iv) sputtering and (v) electric
explosion of wire.
(i) Mechanical Milling (MM) Milling is the process of reducing relatively coarse
materials to desired fineness and is a potential method of producing nanos-
tructures through top-down route. Most of the milling machines employ ball-
milling method for turning bigger lumps of precursor material into nanosize.
A ball media milling machine has a fixed double-wall stainless steel cylinder
fitted with water cooling arrangement. The cylinder is filled with big pieces
of precursor material and large number of stainless steel balls which work as
milling medium. A rotor shaft passes through the centre of the cylinder and may
be rotated by an external motor. Rotating impellers attached to the shaft impart
rotator motion to the large number of metal balls. The shaft may be rotated
by different speeds, depending on the type of the final product. Metallic balls
moving with high speed collide with each other and with precursor material
breaking down big pieces of the material into small pieces. Different type of
forces, like elastic forces, plastic forces, shear forces and chemical forces are
applied to precursor particles by the colliding balls as shown in Fig. 9.19. The
speed of rotation, size of balls temperature of the tank, etc. decide the nature
and size of nano- and microstructures found in powdered precursor.
(ii) Laser Ablation Laser ablation is a method of producing nanostructures with
very high purity. The method may be classified both: a top-down technique or as
a bottom-up. It may produce different types of nanostructures including semi-
conductor nano quantum dots, nano wires, nanotubes and core shell nanopar-
ticles. This method utilises laser as the source of energy for ablating solid
precursor material. An intense and focused laser beam, generally pulsed, is
458 9 Nanomaterials
Fig. 9.19 a Section of ball media milling machine, b generation of elastic and plastic forces by
grinding of balls, c generation of shear force by rotating balls
made to hit the solid precursor material at a point, large amount of energy
from the laser beam is absorbed by a small surface area of the target material
which evaporates. The term ‘ablation’ refers to the removal of surface atoms
and involves not only a single photon process of chemical bond breaking but
also multiple-photon process of thermal evaporation. The ablation process may
be carried out in vacuum, in an inert gas or in a liquid. The characteristics of the
nanostructures produced by this method depend on the purity of the precursor
material, the repetition rate and energy of the laser beam, and the environ-
ment (vacuum, inert gas or liquid) under which the ablation has taken place.
Figure 9.20 shows the basic layout of laser ablation technique and the various
steps of the process.
(iii) Nanolithography and Etching Fundamental idea of fabricating nanostruc-
tures using top-down technique is taken from the techniques used in making
miniature solid state electronic devices. The general method used for making
miniature electronic devices is called Lithography. Nanolithography drives
its name from the Greek words nano (small or dwarf), lithos (rock) and
grapho (to write) that literally means small writing on rocks. The technique
is essentially meant to write on some surface features with dimensions of the
order of nanometres. The writing of nanosize structures may be achieved by
depositing, etching or removing material from some sample. Lithography may
Fig. 9.20 Layout of laser ablation technique
be performed using light (photo lithography), electron beam (e-beam lithog-

raphy), ion beam (I-beam lithography) or X-rays (X-ray lithography). The type
of lithography is decided by the size of the final output product; for example if
the size of the output feature is above 500 nm then one may use photolithog-
raphy, for output size of 200 nm or above X-ray lithography is good. However
for features of size 50–100 nm, one may use either electron or ion beam lithog-
raphy. All types of lithography techniques are implemented in several steps,
for example in case of photo lithography, the first step is to coat the surface of
the desired substrate with a thin polymer layer of either a positive or a negative
resist. Positive resist is that material which is not affected by the incident radia-
tion, in case of photo lithography), the layer of positive resist remains unaffected
by the incident light. A negative resist is a material which is destroyed by the
incident light. The resist layer is covered by a predesigned mask sheet. Mask
sheet is a sheet of some material which absorbs the incident (light) radiations
and does not allow them to pass to the resist. The mask sheet is predesigned
to have transparent features engraved on the sheet through which the incident
radiation (light in case of photo lithography) may pass on to the resist coating.
Transparent features engraved on the mask are created on the resist coating; in
case of positive resist, features engraved on mask are printed on the substrate in
the form of undamaged resist coating and in case of negative resist as absence
of resist material, as shown in Fig. 9.21.
After exposure to light the substrate material with positive or negative image
of the mask is treated with a developer that wash off the unwanted resist mate-
rial. The substrate is then itched to generate the desired micro- or nanofeatures.
As has already been mentioned, either electron or ion beam radiations are used
for developing nanosize feature on the mask. Etching may also be done by
different methods, depending on the requirements. There may be wet chemical
460 9 Nanomaterials
Fig. 9.21 Sequence of photolithography
etching, dry etching using plasma, purely physical using ion beam milling, or
reactive ion etching which is the combination of the two.
(iv) Sputtering Sputtering is essentially a technique of creating nano- and micro-
size particles of a given material and to get them deposited on some surface to
make a very fine nano- or microfilm of the material. Thin films made up of nano
(or micro)-size particles have large number of applications; they are used in
microelectronic industry, solar panels, as oxidation protection films, as antire-
flecting coating on cars, jewellery, mirrors, etc. Further, the process of sput-
tering is also used for identifying materials, for etching, for space weathering,
etc.
A typical sputtering setup is shown in Fig. 9.22. The process of sputtering
is generally carried out in a chamber which is filled by some inert gas like
argon at low pressure. The chamber is provided with two electrodes, the anode
and the cathode. Cathode which is kept at a negative potential with respect to
the anode is also called the target. Nano/microsize particles of target/cathode
Fig. 9.22 Sputtering set-up
material are emitted and are deposited in the form of a uniform layer at the
desired substrate attached to the anode.
In order to start the process of sputtering a high voltage is initially applied
between the anode and the cathode which ionises the argon gas molecules
producing plasma between the electrodes. Plasma that contains positive ions of
argon and electrons glow brightly and is sometimes also called glow discharge.
Positive argon ions in plasma get accelerated towards the cathode because
of its negative potential. Accelerated argon ions hit the target (cathode) and
deposit their energy which initiates the emission of negatively charged nano-
or microsize particles (atoms/molecules of target material) from the cathode.
These negatively charged atoms/molecules of the cathode (target) material
move towards the anode which is at positive potential. On reaching anode,
negatively charged nano- or microparticles get neutralised by giving their extra
electron to the substrate and form a thin uniform film on it.
SAQ: Why argon is filled at low pressure in sputtering chamber?
(v) Electric explosion of wire Electric explosion of conducting wires is a promising

technology for making nanoparticle powders. The technique is based on passing
a high density current pulse through a metal wire which produces excessive
heat in the wire and make it explode. The explosion products pass through a
gas environment producing nanoparticles. The method is not only simple but
possesses several other advantages as well. Firstly, it has high energy efficiency
as very little energy is lost in heating of the surroundings. Secondly, many
experimental parameters, like the strength of current pulse, size and dimensions
of the wire, the environmental gas, etc. which govern the quality, size and other
properties of the created nanoparticles can be easily controlled. Experimentally
it has been found that the spread in the size of nanoparticles produced by this
method is quite small and that the power of nanoparticles is quite stable.
462 9 Nanomaterials
Fig. 9.23 Setup for electric

explosion of wire
In most cases an experimental setup like the one shown in Fig. 9.23 is used,
where sudden discharge of stored electric energy through the wire is made to
make it explode. The set-up may be in open air or may be in an enclosure which
may be filled with a desired gas at a desired pressure.
SAQ: Which method of making nanostructures is most suited for (a) making ultra-
pure nanostructures (b) making powdered nanostructures?
9.4.3 Carbon Nanotubes
(i) Discovery Bonding of carbon atom with itself and with other atoms has
remained a fertile field of research for decades. Three allotropic forms of carbon,
the graphene, the graphite and the diamond were well known for considerable
period of time. However, discovery of C60 , the fullerene by Harry Kroto and
Richard Smally et al. in 1985 in gas-phase carbon clusters obtained in evap-
oration of graphite by intense laser beam in helium atmosphere, gave a new
impetus to the study of other allotropic forms of carbon. Since the yield of
fullerene in laser-induced evaporation of graphite was very small, searches were
conducted to find other means of synthesising fullerene in macroscopic amount.
Later, Wolfgang Kratschmer, Donald Huffman and their co-workers detected
the dominance of fullerene structures in the soot deposited on the walls of the
helium-filled (low pressure) arc chamber when a discharge was passed through
two graphite electrodes. With this simple method it was possible to produce
fullerene in large amounts. With the discovery of fullerene a new allotropic
form of carbon got established. Amorphous and the four allotropic structures
of carbon are shown in Fig. 9.24.
Carbon nanotubes are one of the most important by-products of research on
carbon allotropy. The credit of discovering carbon nanotubes goes to Sumio
Fig. 9.24 Structures of allotropic forms of carbon
Iijima of Japan, who was working as electron microscopist at the NEC labora-
tories of Japan. Iijima got impressed by the technique of producing fullerene in
substantial amount adopted by Kratschmer and Huffman and undertook a project
of detailed study of the soot produced in this method using transmission electron
microscope (TEM). He was motivated to carry out detailed studies as some ten
years earlier he studied soot produced in similar arc discharge between carbon
electrodes and has observed several structures of carbon architecture including
curved, closed nanoparticles and tube like structures. In the initial stages study
of soot taken from the walls of the arc chamber appeared almost completely
amorphous, indicating the absence of any long-range structure. Finding no long-
range microstructures in soot from arc chamber walls, Iijima shifted to the soot
that got collected at the cathode of the discharge chamber in the form of rather
hard rod like lump. Detailed study of this soot from the arc cathode showed
very little amorphous mass but a variety of long-range structures, most striking
of which were long hallow fibers, finer and more perfect than any seen previ-
ously. Iijima announced his discovery of carbon nanotubes in October 1991 in
a meeting at Richmond, Virginia, USA, where he showed beautiful pictures of
these tubes. Amplified photo of nanotubes is shown in Fig. 9.25. A month later
he published a paper in Journal NATURE on his observance of nanotubes in
cathode soot which again attracted scientists to carry out further detailed inves-
tigations of cathode soot which was previously discarded as waste. A coloured
photo of single- and multiple-wall carbon nanotubes is shown in Fig. 9.26.
Method of synthesising carbon nanotubes as adopted by Iijima was not very
efficient and was not able to give substantial yield of nanotubes. Lack of avail-
ability of nanotubes in sufficient amount hampered further research on them.
However, in July 1992 Thomas Ebbesen and Pulickel, two scientists working
464 9 Nanomaterials
Fig. 9.25 Amplified photograph of nanotubes
Fig. 9.26 Single and multiwall carbon nanotube
at the same NEC lab of Japan where Iijima worked, described a method for
synthesising carbon nano tubes in large quantities, of the order of few gram.
They showed that the yield of nanotubes in cathode soot dramatically increases
when the pressure of the helium gas in the arc chamber is increased.
Discovery of nanotubes on one hand opened a totally new and fertile branch
of research that has great application potential, while on the other hand it posed a
big question: why nanotubes (and fullerene) were not discovered earlier when all
facilities and instrumentation which was used in discovering nanotubes were
available for almost last twenty years or so? The answer in case of carbon
nanotubes may be that such tubes were observed earlier also but not much impor-
tance was given to such studies. For example, an author claimed that he observed
thread like carbon structures as a product of chemical reaction between CO and
Fe2 O3 at 450 °C temperature. Methods of producing tube like carbon struc-
tures by some chemical reaction like the one above are called catalytic methods
and were known for long. However, the thread like carbon structures from
catalytic reactions were found to be rough, imperfect as compared to the carbon
nanotubes obtained through fullerene path and therefore do not leant them-
selves to potential applications. It is also reported that Roger Bacon, National
Carbon Company, Cleveland, Ohio, in 1960 synthesised highly perfect graphite
whiskers using the technique which was very similar to the arc discharge tech-
nique of Iijima. It may however, be mentioned that whiskers are very different
from nanotubes, primarily, they are much larger, of the order of 5 µm in diameter
and few centimetres in length, as compared to nanotubes that have diameters in
the range of 2.5–30 nm and lengths from few tens of nm to few micrometer.
In his earlier experiments and investigations on thin carbon films, prepared
using arc evaporation in vacuum (not in low-pressure helium atmosphere) Iijima
(1970–1980) reported the presence of some tube like structures on carbon films,
which were treated as some sort of contamination. In summary it may be said
that carbon nanotubes were observed in some earlier than 1991 but were not
paid much attention.
(ii) Characteristics of Carbon Nanotubes Transmission electron microscope
images of the cathode soot at low magnification show a number of tubes
entangled with each other, accompanied with other material including carbon
nanoparticles, hollow fullerene based structures and some disordered carbon.
At higher resolution, one may see the inside structure of carbon tubes; generally
there are several cylindrical nanotubes of reducing diameter imbedded within
each other. The nanotube length is typically of the order of several µm and
their diameters range from 2.5 to 30 nm. The ends of some of the concentric
cylinder like multi wall nanotube structure are closed by sections of fullerene
molecule. Carbon nanotubes are often referred as molecular carbon fiber and
consist of tiny cylinders of graphene closed at each end with caps which contain
six pentagonal rings. A carbon nanotube may be formed by cutting a fullerene
molecule into two halves and placing these two halves of fullerene as caps at
the two ends of a graphene cylinder.
Depending on the orientation of end caps and graphene cylinder, theoretically
there may be three different structures of carbon nanotubes; known as arm chair,
zig-zag and chiral structures as shown in Fig. 9.27.
A substantial volume of research on nanotubes is directed towards the study
of their electronic properties. Experimental investigation of electronic proper-
ties of nanotubes preceded by theoretical calculations, several research groups
calculated the electron band structure of nanotubes using what is called the
tight-binding model. These calculations indicated that the band structure of
nanotubes is a function of their diameter and structure. Experimental investiga-
tion of electronic properties of nanotubes was very difficult initially, but by 1996
it become possible to experimentally confirm theoretical findings on electronic
properties of nanotubes.
There has been considerable interest in the conductivity of carbon nanotubes
(CNT). It has been found that electrical conductivity of CNT depends on several
factors like its structure, weather it is single walled or multi walled, types and
number of twists, diameter, etc. CNT may be perfectly conducting like metals or
they may be semiconducting depending on their structure, chirality (degree of
466 9 Nanomaterials
Fig. 9.27 Three different structures of single-wall carbon nanotube
twist) and diameter. It has been found that armchair structure of CNT is a better
conductor as compared to other structures of same diameter. A bundle of several
single wall nanotubes makes a nanorope. The resistivity of single walled nano
tube rope is found to be of the order of 10−4 Ω-cm at room temperature, while
it was found that these ropes may sustain current densities of more than 107 A/
cm2 , may be as large as 1013 A/cm2 . Single-walled carbon nanotubes (SWNT)
contain defects which in some cases make the SWNT to behave as a transistor, a
rectifying diode, etc. It is also reported that SWNT transmit electronic signals at
high speed and therefore may be used as an interconnector of different electronic
devices.
Study of mechanical properties of nanotubes was also challenging but
studies carried out using atomic force microscopy and transmission electron
microscopy have indicated that just like their electronic properties, the mechan-
ical properties of nanotubes are very different than the mechanical properties
of graphite, graphene or diamond structures of carbon.
In a graphene sheet a carbon atom is connected to other three neighbouring
carbon atoms with very strong chemical bond, that is why the modulus of
elasticity of graphene is one of the largest of any known material. A single-wall
carbon nanotube is built up of a graphene cylinder and, therefore, SWNT is
expected to be ultimate high-strength fibers. Applying pressure at the tip of a
SWNT may cause it to bend without damaging the tip. This property has made
SWNT as an ideal probe tip for scanning microscopy.
Research has indicated that carbon nanotubes are perhaps the best conductors
of heat. On account of their unique electrical, mechanical and thermal properties,
carbon nanotubes are finding numerous applications in host of devices.
Electron gun, a system that generates an intense electron beam, is an essential
part of most of the display systems. In conventional electron guns, electrons are
produced through thermionic emission and are then focused using an electron
accelerating voltage network. In case of nanotubes, electron emission may take
place by field emission. Since the tip of the nanotube is very sharp, a small
voltage difference between the tip and an electrode may generate large electric
field that may be sufficient for field emission of electrons as well as for their
acceleration. In this way carbon nanotube-based electron gun systems may
become a part of all display systems.
There is great demand in industry of plastic through which current may pass
(conducting plastic). Conductive plastic is made by adding some conducting
material to the plastic matrix. Carbon nanotubes (CNTs) have proven to be an
excellent additive to impart electrical conductivity to plastics. CNTs have a
very high Aspect Ratio, which means that a lower loading or concentration of
CNT is required compared to other conductive additives to achieve the same
conductivity. A lower amount of additive preserves the toughness of the polymer
resin, particularly as low temperatures.
Atom of no other element of the periodic table bonds to itself in an extended
network with the strength of carbon–carbon bond. Special properties of carbon
bonding and molecular perfection of SWNT makes them material with special
electrical, mechanical and thermal properties. The pi-electron donated by each
carbon atom is delocalised and may move around the entire structure. This makes
carbon molecule to be the first molecule that has metallic-type conductivity.
High thermal conductivity of carbon molecule may be attributed to the carbon–
carbon bond that may vibrate at high frequency, giving the molecule an intrinsic
thermal conductivity that is higher than any other material.
SAQ: Compare the thermal conductivity of carbon nanotubes with silicon

nanotubes; what are the main differences?
(iii) Applications of carbon nanotubes There are innumerable applications of

carbon nanotubes; therefore, it is not possible to discuss them or even to list
them all. Some important applications are discussed here.
(a) Electronics Molecular electronic elements are the key elements of
nanotechnology. As the size of electronic elements shrink and the element
density in a single chip increases, it becomes difficult to connect one
element with the other. CNT, because of their metal-like conductivity,
becomes ideal material for interconnections. Their geometry, electric
conductivity and defect-free structure make CNT the most frequently used
interconnecting material. Moreover, CNT with deliberately introduced
defects may itself work as an electronic element.
(b) Energy storage With the growing danger of global warming, petroleum/
natural gas-based energy generating system is now being replaced
by battery-operated systems; petrol/diesel/gas-based vehicles are being
replaced by battery-operated vehicles. All sort of batteries use two elec-
trodes; two important properties of any electrode material are (i) large
surface area for a given mass and (b) good electric conductivity. Carbon
nanotubes have very high surface-to-mass ratio, as high as 1000 m2 per
468 9 Nanomaterials
gram and also have very good electric conductivity; therefore, CNT is
finding great application in battery industry.
Capacitor is an element that may store electric energy. Recent
researches have indicated that CNT have very high reversible capacity
and are, therefore, extensively used in lithium-ion batteries.
Carbon nanotubes are finding applications in fuel cells. Fuel cell uses
the chemical energy of hydrogen or some other fuel to cleanly and effi-
ciently produce electricity. The only products of a fuel cell are electricity,
water and heat. They may use a wide range of fuels and feed stocks and are
capable of providing power for systems as large as a utility power station
or as small as laptop. A fuel cell works like a battery but they do not but
do not rundown or require recharging. They work and produce electricity
so long as the fuel (hydrogen in case of hydrogen fuel cell) is supplied. In
a typical hydrogen fuel cell, two electrodes, anode and cathode, are sand-
wiched around an electrolyte, the fuel hydrogen is fed to one electrode
and the air to the other. A catalyst at anode (to which hydrogen is suppled)
converts hydrogen molecules into electrons and protons. The electrons
flow through an external circuit constituting current while protons migrate
through the electrolyte to the cathode and combine with oxygen to make
H2 O. Since catalytic action at anode depends on the exposed area of the
anode, nanotube anodes which are good conductors of electricity make
fuel cells more efficient.
On account of their high mechanical strength and toughness-to-
weight property, carbon nanotubes are finding applications in composite
components in fuel cells used in transport applications.
(c) Electron emitter On account of their very sharp tips, field emission of
electrons from carbon nanotubes occurs at relatively very low applied
potential. This makes CNT as an important component of electron guns
required in all display systems. CNT is being used in flat-panel displays,
instead of a single electron gun as in traditional cathode ray tube, there is
a separate electron gun of nanotube for each pixel in such flat displays.
Another important characteristic of nanotubes, their capability of
sustaining with high stability very high current densities as large as 1010
A/cm2 or more, makes them a very suitable material for lightning arrester
material.
(d) Material properties Extraordinary thermal conductivity, high capacity of
sustaining large current densities, unparalleled mechanical strength and
extraordinary surface-to-volume (or mass) ratio, and very high aspect ratio
make carbon nanotubes a wonder material. CNT is directly used in many
electrical and electronic applications, for example as heat sinks in high-
density microelectronics where elements readily acquire high tempera-
tures. Nano tubes are also used as addends for changing electrical and
mechanical properties of composites.
(e) Filters Many industries are engaged and have developed water and air
filters based on CNT. It is reported that filters based on nanotubes not only
stop particles as small as viruses but also kill bacteria’s and viruses.
(f) Biomedical applications In view of the excellent chemical stability, rich
polyatomic structure, and high surface area CNTs either absorb or conju-
gate with a wide variety of medically important molecules, like molecules
of drugs, proteins, antibodies, strains of DNA, enzymes, etc., and may
carry them near to the targeted cell. Drug may either be loaded or attached
at the surface of nanostructure which may deliver the drug to the desired
target cell either via the endocytosis pathway or via the insertion and
diffusion pathway.
Intense research activities are in progress to use CNTs for treatment
of cancer, in delivering chemotherapist drugs selectively to cancerous
cell without damaging the healthy cells of the body. A water-soluble
conjugate of single-wall carbon nanotubes with Paclitaxel (PTX) has been
found to be highly effective in suppressing tumour growth compared to
conventional drugs in case of breast cancer.
Another application of CNTs is their use in antimicrobial delivery.
Functionalised CNTs may be used in vaccination procedures. It is also
suggested that CNTs might themselves have antimicrobial activity through
oxidation of the intracellular antioxidants.
SWCNTs have strong optical absorption of ultraviolet to near-infrared
radiations. The absorbed radiation energy is almost entirely converted in
heat. This heat may be used to carry out photothermal therapy and imaging.
Materials used in dentistry when modified with MWCNT showed better
results on account of enhanced fatigue resistance, flexural strength and
resilience.
(g) Others As mentioned, these are only some of the important applications
of CNTs, many other applications including their use in fabrics for making
them stain resistant, dirt and water repellent, more durable and in devel-
oping scratch-resistant paints for coating on vehicle bodies, etc. are all
developing fast. Huge applications of nanotechnology are finding their
place in armament and space research.
SAQ: In your opinion which characteristic of CNTs is most important and which
application of CNTs is most beneficial for humanity?
SA9.1 Explain the top-downand bottom-up approaches for the synthesis of nanos-
tructures and discuss in details one method of Bottom-up approach for the
synthesis of nanostructures.
SA9.2 How one may define nanomaterials and the nanotechnology? List the
important characteristics of nanomaterials and discuss one of them in some
details.
470 9 Nanomaterials
SA9.3 Give a brief account of the discovery of carbon nanotubes (CNTs). List the
types and characteristics of CNTs.
SA9.4 Which types of microscopes are generally used to study various features of
nanostructures? Explain the working of atomic force microscope and point
out where it is used.
SA9.5 What is an optical tweezers? Discuss the origin of scattering and gradient
forces in optical tweezers.
SA9.6 Explain the sol-gel method of fabricating nanostructures?
SA9.7 What is meant by self assembly? Give with suitable examples, different
types of self assemblies.
SA9.8 Give a list of important methods of synthesising nanostructures using top-
down methodology.
SA9.9 With the help of a suitable figure explain plasma arch method of fabricating
nanostructures.
SA9.10 Give a list of the important characteristics of CNTs. Why tips made up of
carbon nanotube is used in electron scanning tunnelling microscopes?
SA9.11 In few lines explain each of the following (a) Aspect ratio (b) Quantum dot
(c) laser ablation.
SA9.12 Write a note on biomedical applications of carbon nanotubes

Note: In some of the following multiple choice questions more than one alternatives
may be correct. All correct alternatives must be marked for a complete answer in
such cases
MC9.1 Special colours displayed by metallic quantum dots that are not shown by
their bulk material, is due to
(a) Large surface to volume ratio (b) Biomedical properties (c) Localised
Surface Plasmon Resonance (d) Self-absorption
ANS: (c)
MC9.2 Enhancement in mechanical properties of microstructures takes place
because of;
(a) High surface to mass ratio (b) Defect free structure (c) Localised
Surface Plasmon Resonance (d) Self-purification property
ANS: (b), (d)
MC9.3 As compared to the bulk material the energy band gap in semiconducting
nanostructures
(a) Increases because of quantum confinement (b) Decreases because of
quantum confinement (c) Increases because of self-purification property
(d) Increases because of defect free structure
ANS: (a)
MC9.4 Figure (MC9.4) below shows the electron state density for
Figure (MC9.4)
(a) Quantum well (b) Quantum surface (c) Quantum dot (d) normal bulk
material
ANS: (c)
MC9.5 Thermal conductivity of carbon nanotubes is
(a) Small and isotropic (b) large but highly anisotropic (c) high and
isotropic (d) small but highly anisotropic
ANS: (b)
MC9.6 If X, Y and Z respectively represent the magnetic moment per atom
of a quantum dot, quantum wire and 2D surface of some magnetic
nanomaterial, then
(a) X > Y > Z (b) X < Y < Z (c) X > Y < Z (d) X < Y > Z
ANS: (a)
MC9.7 Which of the following instruments may be used to study the surface
topology of nanostructures?
(a) Scanning tunnelling microscope (b) Atomic force microscope (c)
optical tweezers (d) Transmission electron microscope
ANS: (a), (b)
MC9.8 Which of the following instruments may be used to study the internal
structure of nanostructures?
(a) Scanning tunnelling microscope (b) Atomic force microscope (c)
optical tweezers (d) Transmission electron microscope
ANS: (d)
MC9.9 In optical tweezers which force pushes the microstructure towards the axis
of the laser beam?
(a) Scattering force (b) Surface tension force (c) Gradient force (d) Atomic
force
ANS: (c)
MC9.10 Which Bottom-up method produces carbon nanostructures in macro
amount?
472 9 Nanomaterials
(a) Chemical vapour deposition (b) Plasma arc method (c) Sputtering (d)
Wire explosion
ANS: (b)
MC9.11 Spontaneous molecular arrangement of the disordered entities of
molecules into ordered nanostructures is
(a) Self-assembly (b) sputtering (c) Laser ablation (d) Sol-gel method
ANS: (a)
MC9.12 Electrodes made from carbon nanotubes have the following properties
(a) Large surface area (b) good electric conductivity (c) low melting point
(d) mechanical strength
ANS: (a), (b), (d)
MC9.13 Carbon nanotubes may sustain current densities of the order of
(a) 1050 A/cm2 (b) 1040 A/cm2 (c) 1010 A/cm2 (d) 10−50 A/cm2
ANS: (c)
MC9.14 Single wall carbon nanotubes strongly absorb electromagnetic radiations
in the range
(a) Violet to green (b) Green to yellow (c) ultraviolet to near infrared (d)
microwave range
ANS: (c)
LA9.1 What are nanomaterials, in what respect they are different from their bulk
materials and why these materials are gaining so much importance?
LA9.2 What are top-down and bottom-up methodologies of fabricating nanomate-
rials? Give a list of important methods adopted to synthesise nanomaterials
using bottom-up approach and discuss pone method in details.
LA9.3 Explain with necessary details the nanolithography technique of making
nanostructures. This technique of nanosynthesis falls under under top-down
or bottom-up approaches.
LA9.4 Give a detailed account of the discovery of carbon nanotubes and their
special properties.
LA9.5 Discuss in details at least three uses of carbon nanotubes.
Chapter 10
Sustainability and Sustainable Energy
Options
Objective
The concept of sustainability and its social, economical and environmental aspects
are presented in this chapter. Emission of greenhouse gases from various sectors,
threat of global warming, causes and its impact on sustainability of life, along with
sustainable energy sources, are discussed in detail in this chapter. It is hoped that after
reading this chapter the reader will become more aware of his/her responsibilities
towards developing a sustainable society and sustainable environment.
10.1 Introduction
The idea and concept of sustainability originated in the deliberations of World

Commission on Environment and Development, constituted by the United Nations
in 1983 under the leadership of Norwegian Prime Minister Mr. G. H. Brundtland.
The idea behind appointing the commission was to analyse and find reasons why, in
spite of efforts to improve standard of living of world population using large-scale
industrialisation, people in many countries were still facing extreme poverty, malnu-
trition and hunger. It appeared that in most countries which have achieved a certain
level of prosperity, economic development took place at the coast of ecological heath.
Dedicated efforts to establish social equity also did not yield the desired results. A
common feeling was that the world is required to find a way to harmonise ecology
with prosperity. The commission submitted its final report after four years, in 1987
with the title ‘Our common future’ and defined sustainable development as:
‘Development that meets the needs of the present without compromising the ability of
future generations to meet their own needs’.
Overall sustainability of life on planet Earth is linked with sustainable develop-

ment of social, economical, environment, energy, food and health sectors in all parts
of the world. It is to be realised that lopsided development of any one of the above and
https://doi.org/10.1007/978-3-031-32084-2_10
474 10 Sustainability and Sustainable Energy Options
only in a certain specific part of the world is bound to lead to a disaster of unimagin-
able magnitude. Division of world into developed, developing, underdeveloped and
non-developed has already created a situation where sustainability of any kind of life
on any part of the globe is seriously threatened; life is at razors edge in every part of
the world, including the so-called developed countries.
10.2 Social Sustainability
Sustainability is not only an issue of environment or economy; it is also a social issue.

Extreme poverty is in itself unsustainable; poor neighbourhoods become breeding
grounds for crime and sickness leading to the destruction of infrastructure. Countries
that pay attention to the social development of their people become more sustainable
than those who do not care for their social environment. In a sustainable society,
everybody must have at least minimal access to safe housing, health, police protection
and reasonably nutritional food. Cities with food deserts, where people cannot reach
a food outlet by walking, in low income areas are unstainable.
Social sustainability defines the longevity and wellbeing of a society. A commu-
nity/society have two components: people and the places where they live. Commu-
nity, therefore, has both social and physical environments both of which need to be
sustainable. Physical sustainability alone may give better choices to the community
for a sustainable life style, but the social system will ultimately collapse if there is no
social sustainability. Social sustainability deals with how individuals, societies and
communities live with each other, within the given land boundaries, and cooperate in
achieving the objectives of the development model they have chosen for themselves.
Essentially, social sustainability focuses on processes for creating sustainable and
successful places that enhance wellbeing by understanding the needs and expec-
tations of the people from the places they live and work. This requires designing
and developing the infrastructure required to support social and cultural life, social
amenities, systems for the useful engagement of people, leaving enough space for
people and places to develop further.
Professor Amartya Sen, Nobel Laureate, has spelt out the following five
dimensions of social sustainability:
(i) Quality of Life Affordable housing, medical support, education and training
facilities, employment, ensuring dignity and safety are some of the important
attributes of quality of life. Other facilities like libraries, playgrounds, eating
joints, etc. with enough green environments add to the quality of life.
(ii) Social Cohesion A society consist of several groups of people having different
ways of living, of religious beliefs, of professions, etc. Cohesion requires
that these different groups understand each other, work in cooperation for the
development of the society, help each other in resolving mutual problems, etc.
(iii) Equality It always happens that some group or groups in a society are in a
state of disadvantage; they do not have full control on their lives. Principle of
10.3 Economical Sustainability 475
social equality requires that other groups, those which are in some advanta-
geous position, should help in removing or lowering the barriers so that the
disadvantage groups may have more control on their lives.
(iv) Diversity Different groups in a society add to its diversity which needs to be
protected. Instead of forcing groups to live exactly in the same way, expec-
tations and needs of all different groups must be ascertained, and infrastruc-
tures for the support and the growth of each group be developed with mutual
cooperation.
(v) Governance With diverse societies and groups it becomes necessary to frame
rules and regulations that promote quality, equality, diversity and social cohe-
sion in the lives of the people and to have an authority (government) to imple-
ment these rules. A democratically elected government which judiciously
implements rules, probes ways to generate revenue from existing resources
without over draining them, takes steps to enhance resources and spend
collected revenue in providing infrastructures for development of all groups
is essential for sustainability of societies.
World is facing large number of issues associated with social sustainability, for
example, the issue of resolving racial biases, economic disparities, difference in reli-
gious beliefs, non-uniform distribution of resources, poverty, etc. Large number of
government sponsored and an equal number of non-government organisations are
making efforts to resolve problems associated with social sustainability. Programmes
like Corporate Social Responsibility, which stress on the responsibility of the corpo-
rate sector towards societies along with profit making, are helping in developing
infrastructure for social sustainability.
SAQ: What is meant by social sustainability?
10.3 Economical Sustainability
Economic sustainability is achieved with low rate of inflation, stable currency and
high level of employment with corruption-free governments. Corrupt system where
economic resources are plundered leads to anarchy the like of which was seen recently
in Sheri Lanka. Earlier, in 2009, foreclosure crisis in USA leads to financial collapse
and many fold shrinking of economies not in USA but over whole of the world. The
economic unsustainability got kicked by lenders encouraging buyers without steady
income to borrow money at unfixed rates of interest. Soon mortgage holders got
behind their payments; lenders then foreclosed, all of them simultaneously. This put
a large number of houses on the market, to be sold at lower price. The cascade effect
resulted in shrinking of larger economy; people lost their jobs, their life time saving,
their business, etc.
Economic sustainability is largely determined by government policies and plan-
ning. Neglect of infrastructure development, large-scale borrowings at unfavourable
terms, unplanned depletion of national resources and corruption all leads to the
collapse of the economy. Pakistan and several Asian and African countries including
Bangladesh and Sri Lanka took huge loans from China on unfavourable terms,
allowed their local and traditional industry to die and mismanaged their resources.
The results are known to all.
Circular economy based on recycling and reuse is expected to contribute sustain-
ability in a big way. Circular economy is a system that seeks to obtain the maximum
value from already extracted resources. This involves keeping materials and prod-
ucts in use as long as possible instead of replacing those using new resources. After
they have been utilised, the same materials and products are recovered and regen-
erated as vital resources. They are then reintroduced into the supply chain. Circular
economy contributes to sustainability in many different ways; for example it reduces
and even eliminates waste. Waste is one of the main threats to a sustainable world;
the world generates around 2.01 billion tonnes of municipal waste every year, about
33% of which, by the most conservative estimates, is not managed in eco-friendly
manner. Unmanaged component of waste ends up in landfills, water bodies and
into atmosphere. Circular economy may reduce or even eliminate the unmanaged
component of waste. The three R (3R) strategy of waste management, reduce, reuse
and recycle, may turn waste into an asset. Circular economy encourages the use of
renewable sources of energy. Though all renewable energy sources may not be in the
category of sustainable energy sources, but these sources definitely reduce pressure
and dependence on fossil- or coal-based sources. Circular economy lays stress on
sustainable consumption of resources. Circular economy advocates both for respon-
sible production and sustainable consumption. This also helps in the preservation
and protection of the environment. Circular economy creates jobs in many different
and new fields thus providing opportunities of better life to a significant section of
the population. Overall circular economy contributes to economic growth which is
essential for sustainability of a nation and the world.
One essential point for any kind of economy which is mostly forgotten is to reduce
consumption. Advertisement industry, which according to Statista collected globally
revenue to the tune of US$ 6410.22 billion in the year 2021 and has shown a growth
of around 11%, is the biggest instrument pleading for increased consumption and
thus working against the goal of sustainable world.
SAQ: What is meant by circular economy?
10.4 Environmental Sustainability
Environmental sustainability, unlike social and economic sustainability which to

some extent are regional or country specific issues, affects whole of the World. It is
not possible to isolate environment of Africa from that of America or that of Asia
from Europe and so on. Environment may be defined as a sum total of all the living
(biotic) and non-living (abiotic) elements and their effects that influence human life.
All elements of environment put together may be called the Nature. It is a well
known fact that human beings feel relaxed, at peace when connected with Nature
10.4 Environmental Sustainability 477
or the environment, walking through forests, going to lakes, beaches are examples
of this. With the unprecedented growth of population and unrestricted exploitation
of natural resources, the environment has suffered extensive damage. As such it is
the moral responsibility of the present world population to prevent further damage
to the environment, to ensure our future generations have healthy places to live,
and minimize further damage to the earth’s bio-diverse ecosystem. According to UN
recommended environmental program, environmental sustainability involves making
life choices that ensure an equal, if not better, way of life to future generations.
For the sake of simplicity, the environment may be divided into three parts: the
atmosphere, the land mass and the water bodies.
10.4.1 Atmosphere
Planet Earth is enveloped from all sides by a gaseous blanket called atmosphere.
Atmosphere, which is retained by the gravitational force of earth, protects earth by
generating pressure that allows liquid water to exist at the surface of earth, shields
earth from ultraviolet rays of the sun and, most importantly, keeps earth warm by
holding a part of heat energy provided by the sun. A common name for the mixture
of gases in the atmosphere is ‘air’. Air contains by molar fraction (i.e. number of
molecules) about 78% nitrogen, 21% oxygen, 1.0% argon, 0.04% carbon dioxide
(CO2 ), 1% water vapours, 0.00017% methane and still smaller quantities of other
gases like nitrous oxide (N2 O), etc.
(i) Greenhouse Effect
It is a common experience that the interior of a car left in sunshine with all its
windows closed becomes very hot. It happens because sunlight (that includes ultra-
violet, visible and infrared radiations) passing through the windows of the car deposits
energy in the form of heat on every object inside the car: steering wheel, seats, dash
board, etc. As a result temperature of all objects inside the car rises. Now, there is a law
of physics which says that any object at a temperature above absolute zero emits elec-
tromagnetic radiations of very long wavelengths called thermal radiations. Thermal
radiations are invisible. Interior of the car is thus filled with thermal radiations. These
thermal radiations remain trapped within the inside of the car as glass windows do
not allow thermal radiations to go out through them. In short what happens is that
all components of sunlight, i.e. ultraviolet, visible and infrared radiations, etc., keep
going into the car through windows but eventually they all are converted into thermal
radiations which remain trapped inside, being not able to come out as window glasses
do not allow them to come out. In this way heat energy gets trapped within car’s inte-
rior. This is called greenhouse effect as the same principle of trapping solar energy
in the form of heat is used in farming where greenhouses are made using transparent
plastic sheets which allow sunlight to enter the greenhouse but does not allow the heat
energy to leave through plastic sheets. Greenhouses are used to cultivate vegetables
and other plants that need higher temperatures for their growth, particularly in cold
places.
In case of earth, the greenhouse effect is the way in which heat is trapped close
to earth’s surface by the envelope of some gases called greenhouse gasses. These
heat trapping gases (which work like the glass windows of a car) can be thought
of as a blanket wrapped around earth keeping the earth at a higher temperature.
Greenhouse gases include carbon dioxide (CO2 ), water vapours (water in gas form
H2 O), methane (CH4 ), nitrous oxide (N2 O, NO), ozone (O3 ) and chlorofluorocarbons
(CFCs). Except CFCs all other gases are natural. Warming effect of carbon dioxide
gas is basically responsible for maintaining the average temperature on earth to a
comfortable 15 °C. Remove CO2 from the atmosphere, and the greenhouse effect
on earth will collapse, plunging earth to a temperature of the order of − 20 °C.
Natural greenhouse effect on earth allows all forms of life, including human beings, to
flourish. During the last century or so, overpopulation, overindustrialisation, overuse
of fossil fuels, excessive burning of coal, etc. have resulted in overproduction of
carbon dioxide gas in particular and other greenhouse gases in the atmosphere. The
level of CO2 in atmosphere is consistently rising for decades, trapping extra amount
of thermal energy raising the average temperature of earth.
(ii) Sources of Greenhouse Gases
Sources of greenhouse gases may be divided into two class: natural and manmade.
Natural sources of CO2 are: (i) out gassing from the ocean, (ii) decomposing vegeta-
tion and other biomass, (iii) venting volcanoes, (iv) naturally occurring wildfires and
(v) belches from ruminant animals. There are also natural sinks of carbon dioxide,
which absorbs or remove it from the atmosphere; they are (i) photosynthesis by plants
on land and in the ocean, (ii) direct absorption by the ocean and (iii) creation of soil
and peat.
Use of fossil fuel for generating energy, heating, in running vehicles, etc. is the
primary source of CO2, methane and nitrous oxide gas emission by human activi-
ties. Other human activities like deforestation and clearing of land for agriculture,
running of heavy industries, transport sector, etc. are some of the human activities
that generate carbon dioxide. Some industries like cement industry generate huge
amount of CO2 .
Figure 10.1 shows the contributions of two main manmade sources of CO2 : fossil
fuel burning plus cement industry and deforestation. The depleting strengths of two
types of natural sinks, land-based and ocean-based, are also shown in the figure. It
may be seen from this figure that there is an addition of nearly 2 × 109 tonnes of
CO2 in the atmosphere per year, adding to the greenhouse effect.
If one goes back by about 800,000 years, back to the period of ice age, the natural
emission sources and natural depletion sources (sinks) of carbon dioxide kept an
immaculate balance keeping the carbon dioxide level in atmosphere between 200
and 275 ppm (parts per million). Lower values corresponding to ice age and higher
values to warm periods called interglacial periods. However, from approximately
400,000 year back the CO2 level in atmosphere started rising very slowly, reaching
the previous highest concentration of 300 ppm around 300,000 year back. Today the
Fig. 10.1 Contributions of

two main manmade sources
of CO2 emission and natural
land and ocean-based sinks
CO2 level in atmosphere is around 410 ppm, and this enormous increase in the CO2
level has taken place in a very short time in geological scale. If excessive CO2 has not
been released in the atmosphere by human activities, it might have taken 1000 years
to reach this level. Figure 10.2 shows the variation in the concentration of CO2 since
ice age till now.
The percentage of methane in atmosphere is very low, yet it is the next important
greenhouse gas after CO2 . The main source of methane is the wetland where methane
(marsh gas) is produced by the anaerobic decay of vegetation. The molar concen-
tration of methane in the atmosphere is rising continuously. The 14.7 ppb (parts per
billion) increase in CH4 concentration observed in 2020 was the largest of the past
four decades. Since 1750 its relative concentration has increased twice as fast as
that of CO2 . Like CO2 , methane also has natural and manmade sources. Wetlands,
termites, cattle, sheep, other ruminants and oceans are the natural sources of methane
Fig. 10.2 Carbon dioxide concentration since ice age till now
Fig. 10.3 Relative

contributions of natural and
manmade sources of
methane and the strength of
natural sink
emission. Hydroxyl group (OH) is a natural sink for methane. Methane reacts with
hydroxyl radical (OH) forming water and carbon dioxide. Agricultural activities,
burning of biomass and waste management are the prime causes of methane emis-
sion through human activities. Methane is also emitted during the production and
transport of coal, by the decay of organic waste in municipal solid waste landfills.
Figure 10.3 shows the relative strengths of the natural and the manmade sources
of methane emission. Strength of natural sink of methane is also shown in the figure.
Though CH4 molecule stays in the atmosphere for a short time being readily oxidised
by (OH) radicals in presence of ultraviolet light, but excessive emission of the gas
by human activities has disturbed the balance between the emission from manmade
sources and its depletion from natural sink, enhancing the greenhouse effect.
Natural source for nitrous oxide (N2 O) includes oxidation of ammonia in the
atmosphere and from nitrogen in soils. Natural sink for N2 O is photolysis to nitrogen
(N2 ) and oxygen (O). Nitrous oxide (N2 O) also reacts with oxygen to make NO which
may enter into a stratospheric ozone-depleting reaction cycle. Manmade source for
the release of N2 O in the atmosphere is the extensive use of fertilisers in agriculture.
Some nitrous oxide is also released by the burning of biomass and by the decay of
livestock manure.
Figure 10.4 shows the relative contributions of natural and manmade sources
of N2 O emission. It may be observed that both sources have nearly the same
contributions, which is nearly compensated by natural sink of the gas.
SAQ: What is greenhouse effect? What causes it?
10.4.2 Mechanism of Trapping Heat by Greenhouse Gases
One may ask the question: what is so special in greenhouse gases that they trap heat?
The answer is simple. According to the quantum mechanical picture, molecules
Fig. 10.4 Relative

contributions of natural and
manmade sources of N2 O
emission and natural sink
of greenhouse gases have discrete energy states that may be excited by thermal
radiations. So whenever a molecule of greenhouse gas is hit by a photon of thermal
radiations it is absorbed and the molecule goes to the excited state. These excited
states are short lived, and therefore, excited molecule reverts back to the ground
state by re-emitting the thermal photon. In this way molecules of the greenhouse
gases absorb thermal radiations and re-emit them, and this goes on and on. Thermal
radiation photons remain in interactions with molecules of greenhouse gases, being
absorbed and re-emitted without being lost. Other gases do have discrete excited
states, but the energies of these states do not match with the energy of thermal
photons; hence molecules of these other gases do not absorb thermal photons.
The efficiency of trapping heat is different for four main greenhouse gases. If the
heat trapping efficiency of CO2 is taken as 1, then the efficiency of methane is around
20, of nitrous oxide 310 and for hydrofluorocarbons around 1000.
Apart of the natural presence of greenhouse gases in the environment, human
activities add a huge amount as given below.
10.4.3 Global Greenhouse Gas Emission by Human Activities
(a) Carbon Dioxide Gas CO2 : Fossil fuel use is the primary source of greenhouse
gas emission, and deforestation, land clearing for agriculture, degradation of
soils, etc. are other sources.
(b) Methane CH4 : Agricultural activities, biomass burning, energy use and waste
management are main sources of methane emission.
(c) Nitrous Oxide N2 O: Use of fertilisers is the primary source; fossil fuel
combustion also releases nitrous oxide.
(d) Fluorinated Gases F-Gases: Includes hydrofluorocarbons (HFCs), per-
fluorocarbons (PFCs) and sulphur hexafluoride; main source of emission is
refrigeration and allied industries.
Fig. 10.5 Greenhouse gas

emission by human activities
Relative amounts of different greenhouse gases generated through human

activities are shown in the following pie chart of Fig. 10.5.
Break-up of the contributions of different economic sectors to greenhouse gases
is shown in the chart of Fig. 10.6.
Economic sector-wise contribution of greenhouse gases is shown in pie chart of
Fig. 10.6. As may be observed in this figure, the maximum contribution of green-
house gases comes from agriculture (24%) and generation of electricity (from coal,
fossil fuel) and heating. A significant contribution is added by transport industry,
particularly the road transport based on fossil fuels.
Out of the three main greenhouse gases, carbon dioxide is the most dangerous
so far as the greenhouse effect is concern. It may be noted that though CH4 has
higher efficiency of heat trapping, yet it is not as dangerous as CO2 . The reason is
that a molecule of CH4 may remain in the atmosphere only for a few days before it
interacts with (OH) radical and is removed. A molecule of carbon dioxide may stay
intact in the atmosphere for a century without undergoing any chemical reaction. It
is for this reason that sustainability of atmosphere demands immediate reduction in
the emission of CO2 from human activities.
Eight countries which generate largest amounts of greenhouses gases are:
(i) China 9877 million metric tonne (MMT)
(ii) USA 4.745 MMT
(iii) India 2310 MMT
10.5 Global Warming 483
Fig. 10.6 Economic

sector-wise distribution of
greenhouse gas emission
(iv) Russia 1640 MMT

(v) Japan 1056 MMT
(vi) Germany 644 MMT
(vii) South Korea 586 MMT
(viii) Iran 583 MMT.
Pie chart of Fig. 10.7 shows the percentage of greenhouse gas emission from some
major polluting countries.
10.5 Global Warming
The net effect of the rapid increase of greenhouse gasses in atmosphere has resulted
in the rise by about 1 °C the average temperature of the earth from the pre-
industrialisation period. This rise in average temperature has, in turn, affected the
climate in different parts of the world. Climate of a place is determined by the patterns
of temperature, wind, atmospheric pressure, humidity and rain over a long period of
time. Different regions in the world, depending on the patterns of above-mentioned
parameters, have different climates like, dry, tropical, cold and moderate. Seasons
and the time when they come and go at a place are determined by the climate of the
place. The type of plants and animal lives that may survive at a place is essentially
determined by the seasons of the area. There is always a delicate balance between
Fig. 10.7 Pie chart showing

greenhouse gas contribution
from some major polluting
countries
the species of plants and animals and the intricate ecosystem of region, which may
be seriously damaged by a little variation in average temperature of earth.
The rise in average temperature of earth (by only 1 °C) is called global warming
and has already affected weather, rain and temperature patterns in parts of the world.
One of the most visible signatures of global warming is ‘hotter days’. Rising sea
level and increased ocean temperatures are melting ice caps and glaciers resulting in
the rise of sea level. Several countries surrounded by the sea and parts of countries
near seashores are under threat of being submerged in the rising ocean. Most of
the extra heat and CO2 in the atmosphere due to enhanced greenhouse effect have
been absorbed by the ocean making it more acidic and warmer. Extreme weather
events like cloud bursts, wild fires, cyclones, draughts and floods have become more
frequent and intense. It is estimated that one out of six species is at risk of extinction
because of climate change due to global warming. It is because the climate change
is so rapid that some species could not adapt themselves with the rapidly changing
environment. Global warming and associated changes in rainfall and temperature
patterns have made it difficult for farmers to graze their livestock and to grow enough
food. Prolong periods of draughts: less amount of rainfall has already created scarcity
of water in some parts of the world.
The main cause of global warming is the accumulation of greenhouse gases in the
atmosphere. In order to maintain the sustainability of atmosphere it is essential (i)
firstly, to reduce the emission of greenhouse gases and (ii) then to remove the excess
of accumulated gases from the atmosphere. Since each individual human being,
animals, manufacturing of goods, transportation, etc. contribute to the emission of
greenhouse gases, it is necessary to change our way of life. The contribution of living
bodies and of different entities/events to the greenhouse gases is often measured in
terms of its carbon footprints.
SAQ: There are many gases like hydrogen, nitrogen, etc. in our atmosphere, but
only CO2 is treated as the gas responsible for global warming, why?
10.5 Global Warming 485
10.5.1 The Carbon Footprint
Carbon footprint, according to the World Health Organization (WHO), is a measure

of the impact that a particular body or an action has on the amount of carbon dioxide
released in the atmosphere. It is measured in terms of CO2 produced in tonnes.
In other words, carbon footprint is the amount of carbon dioxide (CO2 ) emission
associated with all the activities of a person, or any other entity, like country, building,
corporation, etc. It includes direct emission, like those that results from fossil fuel
burning in manufacturing, transportation, heating, etc., plus the emissions required
to produce the electricity associated with goods and services utilised. In addition to
CO2 , emission of other greenhouse gases is also counted in carbon footprint.
The carbon footprint of an individual person may be divided into two compo-
nents: primary and secondary. Primary carbon footprint is that component of total
carbon footprint that may be controlled by the person; carbon footprint generated
by travelling, house warming, food cooking etc. are directly controllable activities.
Remaining carbon footprint due to the consumption of goods and services, over
which the person does not have control is called secondary carbon footprint.
10.5.2 Reducing and Offsetting Carbon Footprints
If we desire to give our next generations a good environment/atmosphere, each of us

must reduce his/her carbon footprint. Higher the carbon footprint, the more waste or
greenhouse gases a person is creating and releasing in the environment. Fortunately,
there are several steps that an individual or the society may take to reduce carbon
footprints and offset carbon emission. Some of these steps are mentioned here;
(i) Living a low-waste life style. A part of any waste is ultimately converted in
greenhouse gases. A life style in which minimum waste is generated is best
suited for sustainable living.
(ii) Travelling by air is one of the worst modes of travel as it adds huge amount of
greenhouse gases and avoids it as far as possible. Bicycling is the most eco-
friendly mode of transport, and it being a good exercise generates negligible
greenhouse gases. Car pooling and using public transport instead of own
vehicle are more eco-friendly modes.
(iii) Eating plant based diet helps in reducing carbon footprint because it requires
less resources and land. It takes less time and energy to cook plant based
food. Moreover, Livestock, like goats, lambs, cow etc. which are used making
non-vegetarian food is responsible for almost 14% of manmade greenhouse
gases, mostly methane which is about 20 times more effective in trapping
heat. Reduction in the population of such livestock will go a long way in
reducing carbon footprints.
(iv) Apart from eating vegetarian food, eating locally produced seasonal vege-
tarian food also helps in reducing carbon footprints.
(v) Reducing food waste is an important way to cutback carbon footprints.

Instead of sending your food waste to municipal garbage which invariably
ends up in landfills, it is better to compost the food waste. Composting
converts biodegradable and compostable food into soil, leaving zero waste.
(vi) All appliances emit greenhouse gases; therefore, use appliances marked as
energy efficient which are designed to have minimal emission.
(vii) Living in nicely insulated houses so that minimum energy is used to heat or
cool the house is an indirect way of reducing carbon footprints
(viii) Carbon offset tax. International projects based on growing large number of
new plants/forests have been undertaken by world agencies. Plants are very
efficient in removing CO2 and adding O2 to the environment. To manage
these projects, a tax called carbon offset tax may be paid by the individual or
a society or a country to offset or reduce their carbon footprints.
10.6 Projections on Average Temperature Rise of 1.5 °C

Above Pre-industrial Levels
Earth’s average temperature has already risen by approximately 1.1 °C above pre-
industrial level and the report drafted by more than 200 scientists from over sixty
countries predicts that the world will reach or even exceed 1.5 °C of warming within
the next 20 years, in spite of the global emission reduction efforts. The report predicts
the following sever effects that may happen due to global warming:
(i) Sever Heat Waves The report predicts that almost 14% of world population
will experience sever heat waves at least once in five years of time.
(ii) Floods and Droughts Many regions worldwide will experience excessive
floods and draughts that will result in food shortages.
(iii) Rise in Sea Level Hot weather spells, and higher average temperature will
melt glaciers, resulting in the rise of sea level submerging coastal cities and
threatening small island nations. Top four countries the population of which
will be adversely affected by sea level rise are: China (5.5 million people),
Vietnam (23.4 million people), Japan (12.8 million people) and India (12.6
million people).
(iv) Arctic Ice Thaws At least once in a century, the Arctic will experience summer
without any sea ice, a phenomenon that has not happened at least in the last
two thousand years.
(v) Changes in Oceans The report suggests that almost 90% of all coral reefs will
be wiped out, oceans will become more acidic, and marine life and particularly
fisheries will be severely affected.
(vi) Loss of Species Many life forms will not be able to adapt to the rapidly
changing climate and will parish.
SAQ: List five negative effects of global warming observed by you.

10.7 United Nation’s Efforts 487
10.7 United Nation’s Efforts
Over the last few decades, governments have taken collective steps to reduce the
emission of greenhouse gases to slow down the rise of average temperature of earth.
Though not intended to discuss the climate change issue, the Montreal Protocol,
1987, was a historical accord on environmental issues and laid the foundation for
further discussion on the subject. The accord, ratified by almost all countries of the
world, required them to stop producing substances that damage the ozone layer, like
chlorofluorocarbons, etc. The accord was so successful that almost 99% of ozone-
depleting substances was eliminated. Later in 2016, an amendment to the Montreal
Protocol, called Kigali Amendment, asked parties to reduce the production of the
greenhouse gases and substances like hydrofluorocarbons that add to greenhouse gas
emission.
Major UN initiative on climate change took place in 1992 as UN Framework
Convention on Climate Change (UNFCCC), ratified by 197 countries including
USA, adopted the landmark accord of establishing an annual forum, named Confer-
ence of the Parties or COP, for discussion between representatives of different
countries at international level on climate change and reduction of greenhouse gases
in the environment. These meetings of COP resulted in the development and imple-
mentation of two agreements called the Kyoto Protocol and the Paris Agreement.
Kyoto Protocol that was adopted in 1997 and came into force in 2005 was signed by
all member countries including USA (signed in 1998). However, USA never ratified
it and later even withdrew its signatures. Kyoto Protocol was the first legally binding
agreement which required that developed countries to reduce their greenhouse gas
emission by an average of 5% below 1990 levels and establish a system to monitor
countries’ progress. The important element of the protocol was that developing coun-
tries, including China and India, the two major greenhouse-producing countries, were
not asked to take any action.
Paris Agreement, signed in 2015, is the most significant agreement on restricting
the emission of greenhouse gases. The agreement is legally binding as well as it is
self binding; it requires all countries to set emission reduction pledges. Governments
were asked to set targets, called Nationally Determined Contributions (NDCs),
with the goals of preventing the global average temperature from rising 2 °C (3.6 °F)
above pre-industrial levels and pursuing efforts to keep it below 1.5 °C (2.7 °F).
It also aims to reach global net-zero emission in the second half of the century.
Global net-zero emission means that amount of greenhouse gases emitted equals
to the amount of these gases being removed from the atmosphere. Global net-zero
is also called climate neutral or carbon neutral. The agreement also provides for
assessing the progress, called stocktaking, every five years. An important component
of the agreement is the support which the developed countries should provide to the
developing countries in the form of capability building, which essentially means that
the developed countries must invest in efforts by developing countries for reduction of
greenhouse gas emission. The USA, who initially signed the agreement, withdrew
during the president ship of Donald Trump; however, President Joe Biden again
joined and signed the agreement within one month of his election. Countries like
Libya, Yemen and Eritrea have not yet signed the treaty.
In a related move United Nations Environmental Programme (UNEP) identified
six sectors with the potential to reduce greenhouse gas emission enough to keep the
global temperature rise below 1.5 °C by 2030. Agriculture and food sector can
cut greenhouse gas emission by 6.7 Gt (Gaga tonne) per year. Buildings, cities and
construction sector, which is expected to add about 12.6 Gt of greenhouse gases by
2030, may reduce this by proper planning by almost 60%. The energy sector can
cut greenhouse gases by 12.5 Gt annually. Transport sector is responsible for about
one-quarter of all greenhouse gases. Sectors emission is expected to double by 2050.
Actions are required at every level, government, private and public. Industry sector
may cut greenhouse gas emission by 7.3 Gt yearly by using passive and renewable
energy-based systems for cooling and heating. Forest and land use, the greenhouse
gas emission may be reduced by 5.3 Gt annually if further deforestation is stopped
and restoration of degraded woodlands is taken up. These actions would also improve
air quality, increase water supplies to cities, enhance food security and boost rural
economy.
10.7.1 Outlook Scenarios: Computer Model-Based Scenarios

Prepared by IEA
International Energy Agency (IEA) with the help of international computer and
modelling experts derived the following four scenarios for global warming.
1. (Stated Policies Scenario, SPS) Scenario based on the policies stated by different
governments in 2021. This will lead to temperature still rising when it hits 2.6 °C
above pre-industrial levels in 2100.
2. (Announced Pledges Scenario, APS) In this modelling it is assumes that
commitments made by different governments will be met in full and on time.
And if so, the average temperature will rise to the value of 2.1 °C by 2100 and
will continue to increase.
3. (Sustainable Development Scenario, SDS) This model is based on the following
assumptions: (i) all countries have successfully implemented their commitments
on time. (ii) Developed countries reach zero CO2 gas emission by 2050, China by
2060 and all other countries latest by 2070. In such case the average temperature
will peak at 1.7 °C by 2050 and could decline to 1.5 °C by 2100. It estimates that
the energy mix in 2100 will be 58% renewable, 8% nuclear and the remaining
34% by other sources.
4. (Net-Zero Emission by 2050 Scenario, NZE) If so, the temperature will peak
at 1.7 °C by 2050 and will decline to 1.4 °C by 2100. Energy mix will be 50%
renewable, particularly about 70% from wind and solar PV, about 20% from
other renewable sources and most of the remaining from nuclear sources. The
other half will be biomass, gas and oil with carbon capture and storage. Use of
10.8 Sustainability of Land Mass 489
coal falls by 90%, oil by 75% and gas by 55%. Emission from transport sector
will fall by 90%.
10.8 Sustainability of Land Mass
The term ‘land mass’ generally includes rocks, soils, minerals, vegetation and animal
habitats. The condition or land health invokes the concept of ecosystem—the inter-
actions and connections between the living and non-living components of the envi-
ronment. In degraded land, where ecosystem has changed, the altered ecosystems
continue to function but have a reduced capacity to supply the goods and services
sought by livestock. In addition to providing physical needs of increasing population,
land also has spiritual and cultural values for the local population.
Land mass is central to addressing sustainability issue, including biodiversity,
climate change, food security, poverty alleviation and energy. The United Nations
defines sustainable land management (SLM) as ‘the use of land resources
including soils, water, animals and plants, for the production of goods to meet the
changing human needs, while simultaneously ensuring the long term produc-
tivity potential of these resources and the maintenance of their environmental
functions’.
The productivity and sustainability of land mass is determined by the interaction
among three components: land resources (soil, water and biodiversity), climate and
human activities. In face of climate change it is important to select the right land
use for the given socio-economic and biophysical conditions so as to minimise the
land degradation, rehabilitating the degraded land ensuring the sustainable use of
land resources and maximising the resilience. Human activities of sustainable land
use and management in principle decide the resilience or sustainability and degra-
dation of land resources. Sustainable land management (SLM) involves established
practices of soil and water conservation, integrated landscape and natural resource
managements. The four pillars of sustainable land management are:
(i) Institutional support for targeted policies, inbuilt incentive mechanism for the
adoption of SLM and income generation at the local level.
(ii) Integrated use of natural resources at farms and at the ecosystem scale.
(iii) Involvement and participation of land owners, technical experts and policy-
makers with multistakeholders at multilevel.
(iv) Holistic land-user-driven and participatory approaches.
Food and Agriculture Organization (FAO) of the United Nations has the mandate
to support its member countries in developing norms, standards and policies; provide
technical advice; and help in capacity development to implement SLM.
Land degradation has many definitions, but all share the idea that detrimental
changes to the condition of land have occurred because of the many ways land has
been developed and used. Frequently changes in the condition of the land are linked
to a reduction in the productive capacity and its economic value. Important causes
of land degradation are: deforestation, excessive use of fertilisers and pesticides,

Stalinisation, desertification, water logging, soil erosion, overgrazing and waste-
land. Land degradation is a global issue for a number of reasons but most significantly
because productive land is one of several resources where a reducing supply threatens
our capacity to feed a growing world population estimated to be over 9 billion by
2050. Some studies have indicated that out of some 1.5 billion hectares of global
farm land, almost 25% is affected by serious degradation, while the figure was only
15% two decades back.
In view of the seriousness of land degradation issue, many countries have included
land repair and slow down of degradation in their legislation and in secondary and
tertiary education.
SAQ: List at least three precautions that should be taken in land use for its
sustainability.
10.9 Sustainability of Water Bodies
Water bodies consisting of oceans, rivers, lakes and pounds are very important for
the survival of life on planet Earth. Unfortunately, all these resources of water have
got degraded by their excessive exploitation and natural calamities. It is essential to
slow down their further degradation and repair as far as possible the damage they
suffered.
Oceans provide and regulate our rainwater, drinking water, weather, climate,
coastlines, substantial amount of sea food and even the air we breathe. Unfortu-
nately, there is continuous deterioration, particularly of coastal waters by pollution.
Moreover, sea water is becoming more acidic having adversary effect on biodi-
versity and ecosystems. Oceans are home to nearly a million known species and
contain vast potential for scientific research and discoveries. Over 3 billion people
depend on marine and coastal diversity for their livelihood. Oceans absorb about
23% of annual CO2 generated by human activities. It also absorbs about 90% of
excess heat produced in the climate system. Excess heat in oceans is causing heat
waves for marine life, destroying coral reefs, threatening the marine life and severely
affecting the ecosystem. According to estimates about 5–12 million metric tonnes
of plastic enters the oceans. This results in a loss of around $13 billons per year in
cleaning operations and fisheries industries. Heavy tourism hubs lie in coastal areas,
and about 80% of total tourism occurs in coastal areas. Coastal tourism industry
grows at a fantastic rate of US$130 billion per annum. Coastal areas support huge
labour force, almost one third of the total labour force. It is, therefore, very essential
that oceans, seas and coastal areas are properly managed. Unplanned/uncontrolled
coastal tourism may destroy natural resources of the area unless measures are taken
now to regulate it.
Oceans are very intimately related to human health on the planet. Diversity of
species found in ocean is a great boon to pharmaceuticals industry. Bacteria’s found
10.9 Sustainability of Water Bodies 491
deep inside the ocean are very important for developing test kits; for example, these
bacteria were used to develop rapid testing kits for COVID-19 cases.
In order to develop a sustainable ocean regime, efforts at two levels are required.
Steps at international level require international cooperation to protect habitats that
are in danger by establishing comprehensive, effective and equitable systems so
as to conserve biodiversity and ensure a sustainable future for marine life. On the
local level, single-use plastic should be strictly banned, sea food consumption be
regulated, a list of certified sea food items be made, and public be encouraged to
buy only certified items. Fishing, particularly of the seed-carrying fish, be controlled
so that the fish population does not deplete. Keeping beaches clean and healthy will
automatically improve the health of the sea/ocean.
10.9.1 Sustainability of River and Other Water Systems
Along with direct rainfall (green water) rivers, wetlands, lakes and aquifers (blue
water) are primary sources of water for human consumption including irrigation.
Rivers and associated wetlands also provide many ecosystem services along with
religious and cultural values. On account of increasing population and increasing
needs of river water for irrigation, rivers are under pressure, and often there is not
enough water and of adequate quality to fulfil the need of the region. This pressure
resulted in the decline of river conditions in many parts of the world.
River sustainability is concerned with resource sufficiency, resilience to water
related risks, access to water supply and other services, the productive use of
water and equitable distribution between different users and generations.
Sustainability of a river or river basin is determined by weather; the river system
can support the long-term ecological and socio-economic functions of the river basin.
Any programme on sustainability of river system must focus on the unique flow
requirements for the river and then create operating plans for dams that achieve
environmental flows, quantity and quality of water flow that must occur downstream
and upstream of dams in order to revive and sustain critical ecological functions
and habitat for species. Since the terrine, the fauna and flora and other parameters
are different for different river systems, in general it is not possible to give one
prescription for the sustainability of all rivers; each river system must be considered
separately, and adequate strategy for arresting further degradation and improving
the degraded system must be evolved. Close collaboration of government agencies,
scientists and other stakeholders is essential along with sustained and long-term
financial support. Many countries have already taken up programmes for improving
river basins, construction of dams, monitoring natural water flow, timing and quantity,
etc. in order to implement river sustainability measures.
10.10 Some Efforts for Improving the Sustainability

of Environment
As predicted by scientists, it is necessary to keep the global temperature rise below

1.5–2.0 °C (above the pre-industrialisation value) to avoid the worst impacts of
climate change, and it is required to first reduce and then to eliminate greenhouse
gas emission completely. However, it is estimated that these efforts alone will not be
able to attain the targets of temperature rise; it will be required to remove some of
the carbon dioxide gas and carbon from the atmosphere. In fact, most of the climate
model scenarios show that billions of metric tonnes of CO2 already present in the
atmosphere should be removed from the atmosphere per year, while ramping up
emission reduction to achieve the target of temperature rise below 2.0 °C above the
pre-industrial level.
Figure 10.8 shows the projections of CO2 emission in the atmosphere as a function
of time. If no efforts of reducing/eliminating of CO2 are made its emission may reach
around 56–57 Gt per year. If different countries take measures to reduce CO2 emission
as pledged by them according to Paris Agreement, huge amount of CO2 emission
will be avoided as shown by the grey coloured part of the curve. Red line in the graph
shows the desired CO2 emission profile if it is desired to reach zero emission targets in
2050. Further, to control global temperature rise below 2.0 °C, considerable amount
of CO2 (shown by greenish blue colour) has to be removed from the atmosphere.
There may be several ways of removing carbon/carbon dioxide from the atmo-
sphere, each of which faces challenges and limitations. The five important methods
often suggested for reducing/eliminating already present CO2 from the atmosphere
are:
(i) Direct Air Capture This method of removing carbon/carbon dioxide from
ambient air is similar to the methods used in capturing carbon/carbon dioxide
at the source like power plants, etc. The technology is there, but it is costly
Fig. 10.8 Projected reduction of greenhouse gases in the atmosphere

10.10 Some Efforts for Improving the Sustainability of Environment 493
and energy-intensive; studies made in 2018 indicated the cost of direct capture
around $100–$240 per metric tonne. It is required to carry out further research
in order to reduce the cost.
(ii) Forests and Farms Trees are particularly excellent at storing carbon dioxide
through photosynthesis. Forests have very high potential of removing CO2
from the atmosphere. Forest cover gets reduced essentially for two reasons:
providing additional land for agriculture and for human habitat. Increasing the
yield of existing farms and using barren lands for human dwellings may reduce
pressure on forests. Planting trees around agriculture farms may enhance CO2
capture.
(iii) Bioenergy with Carbon Capture and Storage (BECCS) This is another way
of using photosynthesis for carbon capture. BECCS is based on utilising the
biomass as a source of energy in industrial, power or transport sectors and
cap the resulting emission of greenhouse gases to reach the environment. The
captured carbon/CO2 may be stored underground or in long-lived products
like concrete.
(iv) Mineralisation of Carbon Some minerals in nature are found to react with
CO2 gas turning carbon of the gas in solid. This process, called carbon miner-
alisation, happens in nature at a very slow rate, and it takes hundred and thou-
sands of years. Scientists are trying to accelerate this slow natural process
by finding suitable catalysts, exposing carbon dioxide containing atmospheric
air to mineral/rock pieces/slag, etc. littered by mining. It has been found that
atmospheric air when injected in some special types of rock structures turns
its carbon into solid carbonates. Since chemical reactions that turn gaseous
carbon into solid take place at the surface of the minerals, nanoparticles of the
minerals that have large surface area are also being investigated for achieving
faster carbon mineralisation.
(v) Capture by Ocean A large part of globe is covered with oceans where different
types of lives thrive. Some of them, like costal plants and weeds, absorb CO2
via photosynthesis. Increasing costal plant density at appropriate places may
enhance carbon capture by a large amount. Similarly, adding some chemicals
to increase storage of dissolved bicarbonates may help in carbon capture.
There are many other ocean-based options which need further research to
make oceans a carbon store without compromising the marine life.
Each carbon capture technique has its own drawbacks and advantages; is suited
to a specific region; and has its own economic viability and the extent to which it can
be used. A mix of all techniques and further research to tailor a particular method to
suit the area are required to achieve the optimum carbon capture.
SAQ: How ocean helps in removing carbon dioxide from the atmosphere?
10.10.1 A Unique Fight Against Climate Change; the Ice

Stupa or Artificial Glacier
There are several places in the world which have zero greenhouse gas emission but
are facing extreme conditions because of the global warming and climate change.
One such place is the area of Ladakh, a high plateau in north India. Climate change
and rise in average global temperature around 1.8 °C above the pre-industrialisation
have dried the natural resource of water in this area. The average annual rainfall in
this region is only around 12 cm, and the flat land is called the desert of north. The
water lifelines of this area were the winter snow and glaciers, which have now almost
disappeared or receded. Some forty years back there use to be enough winter snow
which persists till spring and provided enough water for cultivation of crops. Water
fountains fed by glaciers provided water during summer. However, due to global
warming, winter snow does not stay till spring, and receding glaciers do not provide
sufficient water for summer (Fig. 10.9).
Some youths of the region took the challenge head on. They observed that during
nights and under shade the ice remain frozen in Ladakh, at the height of summer
even at lower altitudes. They thought that water frozen during winters and protected
from sun may be used for irrigation in spring and summer. It was not possible to erect
big shades for the large volumes of ice; however, it was realised that tall mounds of
ice may themselves protect their interior from sun. Simple science that they studied
at their high school level told them that conical shape of the ice mound is best for
shielding its own interior. The idea of conical shape for the ice mound also has some
connection with stupa, conical mounds of stones where religious relics are kept under
the stupa in Buddhist culture. The first ice stupa was built in November 2013. Channel
to bring down water of melting glaciers from higher slopes or rivers was brought
down and was then sent up vertically in a metallic pipe that had a nozzle at the top.
The nozzle was opened during nights when the air was below freezing point. The
fine spray of water coming out of the nozzle froze as it fell. Slowly a mound of ice,
conical in shape, rose around the pipe. The first ice stupa was only 20 feet high, but
Fig. 10.9 Ice stupas in

Ladakh
10.11 Sustainable Energy 495
since then ice stupas or artificial glaciers as high as 100 feet have been made. It is
believed that if ice stupa of optimum height is made at a favourable place, it may last
round the year and accumulation of ice for few years may convert it into a manmade
glacier.
10.11 Sustainable Energy
In line with the definition of sustainability, sustainable energy is the energy that
powers the needs of the present generation maintaining the ability to meet those of the
generations to come. Sustainable energy is replenished constantly via natural means.
In other words sustainable energy will never run out. It is however important to realise
that sustainable energy does not necessarily mean that it is 100% environmentally
safe; it does not pollute environment at all. Sustainable energy definitely has a smaller
carbon footprint as compared to sources like coal or fossil fuels. Sustainable energy
source, which include many renewable energy sources, is in a way the global ticket
to a cleaner and less polluted Earth. Sustainable energy sources are also called the
alternative energy sources.
10.11.1 Units of Energy
Basically, energy is the capacity to perform work, and in science it may be measured
in several units like erg, joule, electron volt (eV), kilowatt hour (kWh), etc. The SI
unit of energy is joule and 1 J = 107 erg = 2.78 × 10−7 kWh = 6.24 × 1018 eV. Some
really big units for energy, like terajoule (TJ)[1TJ = 1012 J], terawatt hour (TWh) {1
TWh = 1012 W h = 109 kWh} and million tonnes of oil equivalent per year (Mtoe),
are also used. 1Mtoe = 11.63 TWh = 11.63 × 109 kWh. Big energy units are often
used in expressing national or international energy production/consumption.
10.11.2 Primary Energy
Let us take the example of crude oil, it is a form of energy, but it is processed to
make it suitable for consumption by end user. The supply chain between production
and final consumption involves many conversion activities and much trade and trans-
port among countries. In energy statistics primary energy refers to the first stage
where energy enters the supply chain before any further conversion or transformation
process.
10.11.3 Global Energy Production, an Overview
Before discussing sustainable energy options, let us have a look on present-day global
energy production and utilisation. According to the International Energy Agency
(IEA) data for year 2018 the total primary energy (PE) production of the world
was 14,420 million tonnes of oil equivalent per year (Mtoe). Data for some other
countries for the same year along with breakup of the contributions from different
energy sources is given in Table 10.1.
As may be observed in Table 10.1, most of the primary energy (PE) till 2018 has
been produced by fossil fuels and coal.
Figure 10.10 shows how the contribution of different sources of primary energy
increased with time and the relative contributions of different source in 2020. As is
evident from this chart almost 83% of total PE in the world is still derived from the
three most polluting fossil fuels: oil, coal and natural gas.
Average per head consumption of energy in some countries and regions is shown
in Table 10.2. As may be seen in this table, energy consumption per head in developed
countries is much higher than that of people in developing countries.
Last two frames of Fig. 10.11 show the history of global energy production: the
landmark events that helped in energy generation to meet the growing demand of
energy ignited by the industrial revolution of 1750.
SAQ: Which source of primary energy is most polluting?
Coming back to the issue of sustainable energy options, sustainable energy
sources must be such that they never exhaust; i.e. the nature must keep sustain-
able energy sources alive. It is obvious that such sustainable energy sources must
be derived from those natural objects and events that are expected to continue undi-
minished for generations to come. Heavenly bodies like sun, moon and earth (all
expected to live for another 3 billion years or more) and events like rotation of earth,
gravitational pull of earth and moon on oceans, etc. are expected to continue for many
future generations of mankind. Energy sources like sunshine, wind and ocean waves,
Table 10.1 PE production in some countries/region in Mtoe

Country Total Coal Oil and gas Nuclear Renewable
China 2560 1860 325 77 300
USA 2127 369 1400 219 180
Middle East 2040 1 2030 2 4
Russia 1484 240 1165 54 25
Africa 1169 157 611 3 397
Europe 1111 171 398 244 296
India 574 289 67 10 208
Australia 412 287 115 0 9
Mexico 159 7 132 4 16
Fig. 10.10 Percentage contribution of different primary energy sources in total global energy
production in 2018
Table 10.2 Average consumption of energy per person in different countries/regions

Country or Consumption Country or Consumption Country or Consumption
region per person (toe) region per person (toe) region per person (toe)
China 1.4 USA 4.4 Europe 2.5
Africa 0.5 India 0.4 Russia 3.0
Australia 3.2 Mexico 1.0 Canada 5.0
geothermal energy, hydroelectric energy and biomass energy, energy from hydrogen;
all come under the category of sustainable energy sources as they are inexhaustible
and are derived from natural objects or processes that are likely to continue forever.
Other renewable energy sources are wood, sun-derived biodiesel and grain-derived
ethanol.
10.11.4 Electricity: The Most Convenient Form of Energy
In science it is said that ‘Energy can neither be created nor it could be destroyed, it
may, however, change from one form to the other form’. Large amount of energy is
scattered in the universe in different forms; in the form of heat in sunshine, heat in
the form of geothermal rise of temperature as one goes underground below earth’s
surface, in form of kinetic energy of wind, sea waves and of river water, in the
form of potential energy of nuclear field and so on. Coal, crude oil, natural gas and
hydrogen gas are also sources of energy which is released when they burn in presence
of oxygen. However, energies in all their natural forms are not directly useable; for
example it is possible to utilise solar energy to some extent in cooking food or for
drying clothes, but most of the other energies in their natural form cannot be used
Fig. 10.11 History of energy generation

directly. Moreover, energy in natural forms is confined to specific locations, while

humans require energy in a convenient form that may be transported from one place
to another. Human efforts of producing energy are simply efforts of converting some
of the energy already present in the nature in some less usable form into a form more
useful and suited to human activities. Energy in electrical form, i.e. electrical potential
difference (voltage) and current (electric charge moving with constant speed), is very
convenient and can be easily transported from one place to another. It is for this reason
that most of the power projects convert energy available in different natural forms to
electrical energy.
10.11.5 Cost of Electricity by Source: Cost Metrics
Leveilized cost of electricity (LCOE) is a metric that attempts to compare costs

of different methods of electricity generation on a consistent basis. LCOE is often
presented as the minimum constant price at which electricity must be sold in order to
break-even over the life time of the project. In approximate terms, LCOE is the net
present value of all costs over the life time of the project divided by an appropriately
discounted total of the energy output from the project over its life time.
Since both the cost and the energy output of the project are estimated over the life
time of the project, several approximations, some of which may be purely subjec-
tive, are made in calculating LCOE. Also, both the cost and the output energy from
a project depend on the location of the project, and therefore the estimated LCOE
for similar projects at two different places may be widely different. Several interna-
tional agencies make their own estimates of LCOE for electricity generation through
different sources, which differ considerably from each other. Average value of LCOE
(in US $) for producing 1MWh of energy using different sources, given by different
agencies, is shown in Fig. 10.12.
SAQ: What are different units of measuring energy and power?
10.11.6 Energy Densities Associated with Prevalent Energy

Sources
Almost all countries of the world are planning to replace most of their polluting
energy-generating systems based on fossil fuel by alternate renewable, cleaner energy
sources. Financial aspect of this planning requires a basis of comparison of different
renewable sources in consistent units. One such unit is LCOE. Another measure
may be the energy density on a Joule per cubic metre or Joule per kilogram basis.
Energy density when measured in Joule per cubic metre basis is called volumetric
energy density and in units of Joule per kilogram, the gravimetric energy density.
Energy density of a fuel or a source of energy is in a way a measure of its efficiency.
Fig. 10.12 Average value of

LCOE for different energy
sources
Knowledge of energy density of a fuel is very important; for example, let us consider
the energy density of a battery; higher the energy density of a battery, longer the
battery can emit charge in relation to its size; high density batteries are useful when
much room is not available for batteries to keep, but lot of energy output is required.
Volumetric energy densities of some sources are given in Table 10.3.
Gravimetric energy density of some fuels is given in Table 10.4.
SAQ: Do you see some relation between the LCOE and the energy density of a
primary energy source?
Table 10.3 Volumetric

Energy source Joules/cubic metre
energy densities of some
sources Natural gas 40 × 106
Food 30 × 106
Gasoline 10 × 109
Oil 45 × 109
Tidal wave 0.5–50
Wind at 5 m/s 10
Geothermal 0.05
Solar 1.5 × 10−6
Nuclear 750 × 1010
10.12 Some Clean and Sustainable Sources 501
Table 10.4 Gravimetric

Fuel Gravimetric energy density in
energy density of some fuels
Mega Joule (MJ)/kg
Natural gas 55
Oil 45
Ethanol 25
Coal 22
Wood 15
Nuclear (U235 ) 3.9 × 106
10.11.7 Problem with Present Energy Mix
As mentioned earlier, at present roughly 83% of global energy consumption is derived

from three fossil fuels: natural gas, oil and coal. Apart from being most polluting,
the known reserves of these fuels are depleting fast; it is expected that world reserves
of natural gas will get exhausted in 53 years, of oil in 50 years and of coal in about
110 years. Secondly, all the three sources of energy put together are responsible for
emitting almost 88% of total CO2 emission in the world. CO2 is the most dangerous
greenhouse gas that survives for a very long time in the environment causing global
warming.
Figure 10.13 shows the per capita CO2 emission (in tonnes per year) for six most
polluting countries. It may be observed that though per head emission in Canada
is the highest (≈18.58 tonnes per year), but total emission by the country is only
around 0.6 × 109 tonnes because of low population. On the other hand for India,
where per head emission is only 1.9 tonnes but because of its large population total
CO2 emission by India is very large around 2.5 × 109 tonnes per year.
In order to reduce CO2 emission it is required to replace existing fossil fuel energy
sources by cleaner and sustainable sources.
SAQ: Which industry produces maximum greenhouse gases?
10.12 Some Clean and Sustainable Sources
Though in common perception only renewable energy sources can be sustainable;

however, the fact is that some energy sources, like hydrogen and nuclear sources of
energy which are not renewable, are also included in the list of sustainable sources.
It is because of two reasons: firstly there are enough reserves of the basic fuels in
nature to last for a few generations, and secondly, the gravimetric energy densities
associated with these sources are very high; only a small amount of these materials
may produce huge amount of energy.
Fig. 10.13 Per head CO2 emission in six most polluting countries
10.12.1 Hydrogen as an Alternative Source of Energy
Chemical reaction between hydrogen and oxygen (gases) is an exothermic reaction

resulting in the release of energy in the form of heat and production of pure water.
Since no polluting gas is emitted in the process, burning of hydrogen in the presence
of oxygen is a very clean energy source. Moreover the end product water is a very
stable compound. Hydrogen is just like any other fuel; for example it is like petrol
which can be transported from one place to another; therefore it is not only a source
of energy but also an energy carrier or transporter. Hydrogen being the lightest
element has a high gravimetric energy density (1 kg H2 produces as much energy
as is produced by 2.8 kg of petrol). But its volumetric energy density is low, and
therefore it has to be stored and transported at high pressures. Hydrogen liquefies at
a temperature of − 423 °C and may be transported from one place to the other either
as gas in pressurised tankers or as liquid in cryogenic freezer tanks. Hydrogen may
also be transported through pipes to places where consumption of the gas is large
and demand remains steady for long periods of time.
(i) Availability: The average hydrogen content of air at ground level is only about
0.6 parts per million, but large amount of hydrogen is present on earth in chem-
ical combination with other elements like in the molecules of water (H2 O) and
hydrocarbons (Cx Hy ). With about 70% of earth being water and almost inex-
haustible possibilities of generating hydrocarbon gases by the decomposition
of biomass and waste, hydrogen reserves may be treated as inexhaustible.
(ii) Production: The two main methods of producing hydrogen are (a) steam
methane reforming and (b) electrolysis of water. Steam methane reforming
10.13 Hydrogen Fuel Cell 503
is the widely used method of producing hydrogen on commercial scale, mostly

in oil refineries. In this process high-temperature (700–1000 °C) and high-
pressure (3–5 bar; 1 psi = 14.5 bar) steam are made to react with methane
(CH4 ) in presence of some catalyst. Chemical reaction between steam and
methane gives hydrogen gas along with CO and CO2 gases in small amount.
In most cases natural gas is used for methane. Methane obtained from landfills
and biogas plants is also used for producing hydrogen.
Electrolysis of water using electric current is also a method of producing
hydrogen, but it is not cost-effective; generally larger amount of electrical
energy gets used up than the energy value of the hydrogen produced in the
process. However, the advantage is that no dirty or greenhouse gas like CO or
CO2 is produced in electrolysis. Electrolysis using nuclear energy is also being
developed.
Hydrogen produced without the emission of greenhouse gases is often
referred as green hydrogen, while hydrogen produced with the emission of
greenhouse gases is called grey hydrogen.
Several other methods based on the use of microbes and light to produce
hydrogen, using solar energy to split water to get hydrogen, etc., are under
research all over the world.
SAQ: What is the reason of transporting hydrogen in high pressure or freezer tanks?
(iii) Fuel cellIt is a device that converts chemical energy released in a reaction
directly into electrical energy, similar to that of an electrical battery or cell.
However, the conversion efficiency of fuel cell is much higher (> 60%) than
that of a combustion engine used in automobiles. Further, unlike ordinary cell,
fuel cell continues to deliver electrical energy at a constant rate so long as
the fuel supply is available. In contrast an ordinary battery supplies electrical
energy at a rate that decreases with time. Fuel cell has no moving parts and is
therefore noiseless. Most of the fuel cells do not produce any greenhouse gas
and other polluting substance and are therefore very clean source of energy. In
principle many different types of fuels may be used to run a fuel cell. Fuel cells
have great potential as they may supply electrical power to a system as small
as a laptop or to a system as big as an automobile, lift fork, air-conditioning
systems at airports etc.
10.13 Hydrogen Fuel Cell
The physical size of a single hydrogen fuel cell is extremely small, of the order of
3 × 3 × 1 mm. Each hydrogen fuel cell (HFC) typically delivers between 0.5 and
0.8 V in its output, and therefore many HFCs are connected in series to get desired
output. The series and parallel combinations of HFCs make stacks which are used to
run devices. HFCs use hydrogen and oxygen (either pure or as air) gases as fuel; the
chemical reaction in the cell produces electrical energy, some heat and pure water as
the by-product.
Fuel cell technology is fast evolving, and new designs and structures of fuel cells
appear almost on daily basis. Basic structure of a hydrogen fuel cell is shown in
Fig. 10.14. A hydrogen fuel cell has two electrodes, an anode and a cathode. Each
electrode is coated with a layer of catalyst which is followed by a layer of polymer
membrane. An electrolyte is sandwiched between the two membranes of the two
electrodes. Fuel cells are classified primarily by the kind of electrolyte used in it.
This classification determines the type of the electro-chemical reaction that takes
place in the cell, the type of catalyst required, the appropriate temperature at which
the cell operates and other factors.
(a) Polymer Electrolyte Membrane (PEM) Hydrogen Fuel Cell

In polymer electrolyte membrane (PEM) fuel cell the polymer membrane
also works as the ion diffusion medium or electrolyte. Hydrogen gas enters in
the fuel cell at anode and oxygen gas (or air) at the cathode. The by-product
water is produced at the cathode from where it is drained out. Both electrodes
of a PEM cell are made up of porous carbon. The anode is coated with a special
catalyst. In presence of the catalyst, an atom of hydrogen got split into posi-
tively charged hydrogen ion H+ (proton) and electron (e− ) at anode. There is a
polymer membrane at the anode on the electrolyte side which allows only H+
ions to pass through it towards the cathode. Electrons formed by the splitting of
hydrogen atom at anode could not pass through the membrane and flow through
the external circuit in the form of electric current. Positive H+ ions move through
the polymer membrane which also works as electrolyte and reach cathode where
they combine, in presence of the catalyst, with the electrons that have reached
cathode through the external circuit and the oxygen atoms already present there
to produce water. Working of a PEM hydrogen fuel cell is shown in Fig. 10.15.
The two electrodes (anode and cathode), catalysts at both electrodes, the
electrolyte membrane (all so-called the diffusion medium as positive hydrogen
ions diffuse through it) put together are called membrane electrode assembly
(MEA). The polymer electrolyte membrane (PEM) looks like a kitchen plastic
wrap, allows only H+ ions to diffuse towards the cathode and blocks electrons.
Fig. 10.14 Basic structure

of a hydrogen fuel cell
10.13 Hydrogen Fuel Cell 505
Fig. 10.15 Structure of a

hydrogen fuel cell
H+ ions migrate towards the cathode because of the difference in their concen-
trations at anode and cathode, i.e. by diffusion. PEM-type HFCs used in auto-
mobiles have polymer membranes of thickness of the order of 20 µm. A layer of
catalyst is added on both sides of the membrane: the anode layer on anode side
and the cathode layer on cathode side of the membrane. Conventional catalyst
layer is made by dispersing nanosized particles of platinum on a high surface
area carbon support and then mixing an ion-conducting polymer. PEM fuel cells
operate at temperatures around 80 °C or 176 °F.
(b) Alkaline Fuel Cells (AFCs) These fuel cells use a solution of potassium
hydroxide (KOH) in water as the electrolyte for the diffusion of protons (H+ )
from anode to cathode. The advantage of AFC over PEM is that catalysts made
of non-precious metals that are relatively much cheap, as compared to platinum,
may be used in these cells. AFCs were used in American space missions for
providing electric power and water on-board. One big challenge to this tech-
nology is the severe poising of the cell action by the presence of a very little
amount of CO2 in the system. It is, therefore, very important to check that both
the hydrogen and the oxygen or air used in the cell are totally free from any
traces of carbon dioxide.
(c) Phosphoric Acid Fuel Cells (PAFCs) It is one of the most matured types of
the fuel cell and is often referred as the first generation of modern fuel cells.
PAFC uses liquid phosphoric acid as electrolyte. The acid is contained in a
Teflon-bonded Silicon carbide matrix. Porous carbon electrodes and platinum-
based catalysts are employed in these fuel cells. These cells are more tolerant
to impurities and may have efficiency of the order of 80%. However, their
gravimetric energy density is low as a result stack of these cells is large and
heavy. PAFCs are, therefore, used either in big automobiles like trucks/buses or
are used in powering stationary systems like buildings, etc.
(d) Molten Carbonate Fuel Cells (MCFCs) These are temperature fuel cells that
use an electrolyte made of molten carbonate salt mixture that is suspended in a
porous and chemically inert ceramic-lithium-Aluminium oxide matrix. MCFCs
operate at temperatures of the order of 600–700 °C range and use catalysts made
of non-precious metals. They have high efficiency, if the waste heat produced
in the electro-chemical reaction is captured; the efficiency of MCFCs reaches
approximately 85%.
Another big advantage of this fuel cell is that natural gas or biogas (instead
of hydrogen) may be directly used in the cell or generation of electrical energy.
This becomes possible because of the high temperature of operation; at such
high temperature natural or other gases undergo internal reforming generating
hydrogen with in the fuel cell.
Durability of MCFCs is substantially compromised because of the use of
molten carbonate which is corrosive and damages the components of the cell.
(e) Direct Methanol Fuel Cells (DMFCs) DMFCs are powered by methanol
(instead of hydrogen) without it being converted into hydrogen. Methanol is
liquid, has higher energy density and is easy to transport using the existing
systems of transport. DMFCs are used to run small equipment like laptops, etc.
(f) Solid Oxide Fuel Cells (SOFCs) They are fuel cells that operate at very high
temperatures of the order of 1000 °C, eliminating the use of precious metal cata-
lysts and facilitating the use of any hydrogen-rich gas (natural gas, biogas, coal
gas, etc.) as fuel. SOFC uses a hard, non-porous ceramic compound as the elec-
trolyte. They are the most sulphur-resistant fuel cells and are not damaged by the
presence of carbon monoxide. High-temperature operation, however, increases
the start-up time and requires good thermal shielding to protect people and retain
heat. High operating temperature also puts stringent conditions on the materials
used in the cell and adversely affects the durability of the cell components.
Efforts are on to develop lower temperature version of solid electrolyte fuel
cells.
(g) Reversible Fuel Cells Normal hydrogen fuel cell uses hydrogen and oxygen to
produce electricity, heat and water. The reversible fuel cells also have the capa-
bility of operating in the reverse order; that is if they are given water and elec-
tricity they electrolyse water and produce hydrogen and oxygen. This reverse
operation capability makes these cells very useful, since at times when there
is excess production of electric power by some renewable energy source (like
solar cell at peak hours or by wind mill at times of high-speed winds, etc.), the
excess electricity may be supplied to the reversible fuel cell which will generate
hydrogen that may be later used for forward operation.
10.14 Nuclear Energy
Considerable amount of energy is released when (i) some big or heavy atomic nucleus
undergoes fission; that is it splits into two nearly equal nuclei, and also (ii) when two
light nuclei undergo fusion; that is the two light nuclei fuse together to make a bigger
nucleus. The energy released in both the fission and the fusion is totally clean; no
greenhouse gas is released in the atmosphere. Though energy released per fusion is
10.14 Nuclear Energy 507
larger and no radioactive waste is produced in fusion, it has not been possible so
far to extract fusion energy in a controlled way on a commercial scale. In contrast,
systems called nuclear reactors have been developed to trap in a controlled manner
the energy released in fission of heavy nuclei; fission energy is with us for almost
60 years or so, with very few accidents. Fission in some heavy stable or (almost
stable) nuclei may be initiated by neutrons and other nuclear particles. Neutrons of
very low kinetic energies ≈ 0.025 eV may enter a heavy (positively charged) nucleus
with ease, as it has no electric charge, and may initiate nuclear fission of the nucleus.
On the other hand, some other heavy nuclei undergo fission when they are hit by
high-energy (≈ 1 meV or more) neutrons. Low-energy neutrons having energies of
the order of the energy of atoms/molecules of gases present in the atmosphere at room
temperature and pressure (≈ 0.25 eV) are called thermal neutrons and are present
in the atmosphere in thermal equilibrium, colliding with the atoms of atmospheric
gases. Those heavy elements the nuclei of which may undergo fission by thermal
neutrons are called fissile atoms or materials. 235 U, 233 U (bred by 232 Th by neutron
capture), 239 Pu and 241 Pu (bred from 240 Pu by direct neutron capture) are fissile
nuclei which may undergo fission when hit by thermal neutrons. Out of the above
four fissile nuclei, 235 U is found in nature in very small percentage; natural uranium
ore contains only 0.7% of 235 U isotope. Other fissile nuclei may be produced using
specific nuclear reactions.
Uranium occurs naturally in low concentration in soil, rocks and in sea water
and is commercially extracted from uranium bearing minerals like uranite. Natural
uranium contains two isotopes: 235 U 0.711% and 238 U 910.284%. Since 235 U isotope
undergoes fission by thermal neutrons, natural uranium must be enriched in 235 U
isotope to use it as the fuel for thermal neutron reactor. Most of thermal neutron
reactors use uranium that has 3–5% enrichment of 235 U isotope.
235
U nucleus when hit by thermal neutrons may undergo fission in ≈ 87% cases,
and in the remaining 13% cases, it may capture the neutron emitting a gamma ray.
In case of fission, the fission products may be different, as shown in the following
set of equations:
It may thus be observed that as a result of the fission of 235 U by thermal neutrons
two heavier nuclei like Ba and Kr or Te and Zr, etc., called fission fragments (FF),
two to three neutrons and considerable amount of energy are released. Almost 95%
of the energy that is released in fission is in the form of the kinetic energy of fission
fragments. When these highly excited fission fragments moving with high speed
collide with the container walls, they release their kinetic energy in the form of heat,
as such most of the (about 200 MeV = 3.2 × 10−11 J) kinetic energy is ultimately
converted into heat. Two to three (average 2.5) fast neutrons of average energy ≈
2 meV are also produced in the fission of each 235 U nucleus. Fission fragments are
highly excited neutron-rich nuclei, and they de-excite by emitting neutrons, but these
neutrons are emitted after a time delay, sometime after the emission of fast neutrons
that are produced at the instant of fission. Thus, two types of neutrons, prompt
neutrons produced at the instant of fission and delayed neutrons emitted by fission
fragments, are generated as a result of fission. Prompt neutrons make it possible to
establish fission chain reaction, while delayed neutrons play a key role in controlling
the fission reaction rate in the reactor.
Nuclear fuel used in thermal neutron reactors is uranium with 3–5% enrichment of
235
U isotope which means that in about 100 atoms of the fuel there will be about 95–
97 atoms of 338 U isotopes and 5–3 atoms of 235 U isotope. In the space where the fuel is
placed in the reactor there are always some stray thermal neutrons. These background
thermal neutrons may interact with the atoms of both the 235 U and 238 U isotopes.
Interaction of thermal neutrons with 235 U atoms results in their fission producing
energy, on an average 2.5 fast neutrons, some delayed neutrons and radioactive
waste. Interaction of neutrons with 238 U atoms of the fuel produces weapon-grade
fissile material 239 Pu and long-life radioactive waste.
Figure 10.16 shows the results of neutron interactions with fuel atoms.
In case the fuel interacts with only stray thermal neutrons, those present in the
surroundings, very few fission events per unit time will take place on account of the
very low percentage of 235 U atoms and very low density of stray thermal neutrons.
However, if fast neutrons produced in initial fission events are made to lose their
Fig. 10.16 a Thermal neutron fission of 235 U nucleus, b decay of fast-moving fission fragments
into radioactive waste, c absorption of neutron by 238U resulting in the production of weapon-grade
fissile 239Pu and other long-life radioactive waste
energy and are converted into thermal neutrons then the rate of fission events may
be increased such that a self-sustained chain of fission reactions may be maintained.
SAQ: Which fission product is crucial in sustaining of fission chain reaction?
Fast neutrons may be rapidly converted into thermal neutrons if they are made to
collide with atoms of those materials which have nearly same mass. It is because in a
collision with another particle of same mass, the neutron will lose maximum energy
in each collision. Materials that may rapidly and without absorption reduce the
energy of fast neutrons are called moderators. Heavy water (D2 O), normal water
(H2 O) and graphite (C) are the materials that are frequently used as moderators.
Figure 10.17 shows how a self-sustained fission chain reaction may be established
in a fuel which is distributed with moderators around it. An important consideration
for chain reaction is the economy of neutrons; the fuel-moderator assembly should
be such that no neutrons may leak it. In order to derive energy from a nuclear reactor
at a constant power level (nearly same number of fission event per unit time) it is
essential that number of fission events per unit time must remain nearly constant.
This brings in the criticality criteria in consideration. Let us assume that at some
instant ‘t’ A-number of fission events have taken place. These A fission events will,
on an average, produce (2.5A) fast neutrons. These fast neutrons will collide with
moderator and may produce say, N 1 (N 1 < 2.5A) thermal neutrons. These N 1 thermal
neutrons are the first-generation thermal neutrons at time ‘t’. N 1 thermal neutrons
will initiate fission in many more fuel nuclei again producing fast and then thermal
neutrons. Let the number of thermal neutrons produced at this stage is say, N 2 . N 2
is the second-generation thermal neutrons. The ratio η
No. of second generation thermal neutrons N2

η= =
No. of first generation thermal neutrons N1
Is called the criticality parameter; if η is greater than 1, then more thermal

neutrons will produce in the next generation, and hence the fission chain reaction
will grow, and rate of fission will increase. In this case the system is said to be in
supercritical state. When η = 1, the system is critical, fission rate will remain more
or less constant, and the reactor will deliver power at a constant level. Finally, if
η < 1 the system is sub-critical; ultimately the chain reaction and power generation
from the reactor will stop. For proper working of the reactor the criticality parameter
should be maintained around 1. In actual situation when moderators are put in the
system and neutrons are not allowed to leak, the system tends to become supercritical.
A supercritical system behaves like an atom bomb and will eventually explode. To
avoid this, arrestor or control rods made up of materials that absorb thermal neutrons
are periodically inserted in the system for some duration of time to keep η around
1. Materials like boron, cadmium, Silver, hafnium and indium are good absorbers of
thermal neutrons and are used in control or arrestor rods.
Figure 10.18 shows the rough sketch of a thermal neutron reactor where water has
been used both as a moderator and also as a heat exchange element. Cold water enters
the reactor vessel and moderates the fast neutrons produced in the fission of 235 U
Fig. 10.17 Building up of a self-sustained fission chain reaction
nuclei in the fuel elements by the stray neutrons. Since fission of each 235 U nucleus
produces two to three fast neutrons which are readily thermalised by moderator water,
the number of thermal neutrons increases almost exponentially initiating fission of
many more 235 U nuclei and setting a self-sustained fission chain reaction. The power
level (fission rate) of the reactor increases with time, and once it reaches the desired
level, arrestor or control rods are inserted in the reactor core by remote operation.
Control rods absorb thermal neutrons leading to the decrease in the fission rate. When
fission rate starts declining below the desired level, control rods are withdrawn and
are re-inserted when fission rate increases above the desired level. Remote-operated
insertion and withdrawal of control rods maintains the power output of the reactor
to the desired level.
In case of any emergency, the reactor may be switched off by inserting the control
roads to the full extent and removing moderator, water in the present case, from the
system. However, both these operations take time of the order of few ten of seconds
to a few minutes because control roads are bulky and removing water also takes time.
Large switch-off time is a drawback of the reactor.
10.14.1 Drawbacks of Fission Reactor
1. Considerable amount of radioactive waste having long-life radioactive elements

is produced. Since 1950 a stockpile of nearly 250,000 tonnes of highly radioactive
waste has got accumulated by the reactors operating in the world. Safe disposal
of radioactive waste is crucial.
Fig. 10.18 Sketch of a water moderator, water-cooled thermal reactor
2. Weapon-grade fissile material like 239 Pu is produced as a by-product. Large

amount of weapon-grade fissile material produced as a by-product of nuclear
reactors has already accumulated in the world.
3. Nuclear power plants are expensive and take considerable time for commis-
sioning; however, they are easy to run and in the long run produce cheap
energy.
4. Comparatively long switch-off time may be a big disadvantage as during some
emergency if it is required to switch off the reactor, considerable damage
including core-meltdown may occur while the reactor is switched off.
10.14.2 Plus Points of Fission Reactor
1. No greenhouse gas is generated or released during the energy generation from a

nuclear reactor. However, weak carbon footprints are left in the construction of the
buildings and other infrastructure of the nuclear plant. This may be compensated
by developing small forests around the reactor building.
2. Chances of leak of radioactive materials including gases are negligibly small, and
it is because nuclear scientists subscribe to concept of concentrate and confine,
unlike coal-fired energy production which uses the concept of disperse and dilute.
There have been three or four major mishaps in nuclear energy industry including
the 1986 Chernobyl disaster. Investigations of the mishap reviled that the design
of control rods of the reactor was faulty and the overall operation of the nuclear
Table 10.5 Energy density

One kilogram of the fuel Approximate amount of electrical
of some fuels
energy in kWh
Fire wood 1
Fossil fuels 3–5
Natural uranium 50 × 103
Plutonium 60 × 105
plant was poorly managed. Long switch-off time also played a negative role, core-
meltdown occurred during the switch-off, and big explosion took place. Another
reason for the leakage of radioactive gases and solids, etc., was the absence of
any concrete containing walls around the reactor. A big nuclear mishap happened
in 2011 at Fukushima nuclear plant, but it occurred because of the tsunami wave
of unprecedented height. Direct loss of human life occurred only in Chernobyl
mishap.
3. Fuel economy is a prime consideration in choosing a source of energy. Energy
density of nuclear fuel is among the highest as is indicated in Table 10.5;
4. Land required for a power plant, particularly in countries having high density
of population, also becomes an important consideration. Land area required for
nuclear plant is comparable to the area required by a fossil fuel power plant of
same capacity. A nuclear power plant of 1000 megawatt capacity may typically
require a land area of about 2 km2 , while a wind farm of same capacity requires
almost 360 times this area and solar photovoltaic plants of same capacity about
75 times the area.
5. Nuclear reactors are working in the world for almost 70 years; the nuclear
technology is time tested and reliable.
10.14.3 Accelerator-Driven Energy Amplifier
Two main drawbacks of a fission nuclear power plant are: (a) production of large
amount of highly active and long-life radioactive waste and weapon-grade fissile
material, proper and safe disposal of which is a big issue and (b) possibility of
radioactive fallout if control and or cooling systems fail. It is however, important
to note that all radioactive nuclei of long half-lives that are generated during the
nuclear fission in a reactor may be made to undergo fission and produce energy if
hit by neutrons of appropriate energy. Also, long-life radioactive nuclei on hitting
by fast neutrons may break into stable or short-life radionuclides. In thermal neutron
reactors, the criticality factor is kept around 1; therefore, there are no extra neutrons
to burn the waste.
To overcome the problems associated with conventional thermal neutron reactors
a new concept is taking shape in terms of the accelerator-based energy amplifiers.
A proton beam of about 1 mA and 800 MeV from an accelerator is made to hit a
heavy target like lead. Interaction of proton beam with lead produces high flux of
spallation neutrons that have an energy distribution. The spallation neutrons energy
distribution is further modified from the scattering by lead block such that there
develop zones around the interaction point which have neutrons of different energies.
If heavy nuclei are charged in zones of appropriate neutron energy, they may undergo
fission releasing energy. The system will not only burn out its own radioactive waste,
but radioactive waste from other conventional reactors may also be charged in the
system to produce energy. Weapon-grade fissile material accumulated in the world
may also be used as a fuel in this system to generate energy. It may be noted that
the neutron economy of the system is not governed by the fast neutrons emitted
in fission, as is the case with conventional fission reactors; instead it depends on
the current and energy of incident proton beam. The system remains sub-critical in
absence of proton beam and becomes supercritical when proton beam is switched
on. As such the fission process going on when proton beam is on may be switched
off instantly if proton beam is stopped. This avoids the possibilities of accidents that
may happen in conventional reactors because of large switch-off time. It is true that
initially a high current high-energy proton accelerator will need considerable amount
of energy to run it, but once fission energy is available it will be much larger than the
energy consumed by the accelerator, and a part of the fission output energy may be
used to operate the accelerator. The basic concept of the energy amplifier is shown
in Fig. 10.19.
Figure 10.20 shows an artist’s view of an accelerator-driven energy amplifier.
Though nuclear energy is not a renewable source of energy but the present reserves
of fissile and fissionable sources in the world are enough to last for almost a century.
Moreover, many countries are developing techniques to extract uranium from water;
India for example claims to have already developed methods to extract uranium from
water with almost 90% efficiency. The safety track record of nuclear energy is quite
impressive; less than 50 people have directly lost their lives during the history of
nuclear energy, maximum of 14 people in Fukushima, Japan, in 1911 accident. In
comparison much more lives have been lost in accidents of coal-run energy industry
Fig. 10.19 Layout of an accelerator-driven energy amplifier

Fig. 10.20 Artist’s view of an accelerator-driven energy amplifier
and hydropower industry in the world. Nuclear energy is one of the cleanest energies
and requires comparatively very small amount of fuel and space; technology is time
tested and may supply uninterrupted, clean power to industry on a long-time basis.
Expected future developments, like accelerator-based energy amplifiers, will make
fission nuclear energy a very clean, safe and sustainable option.
SAQ: How accelerator-driven amplifiers will make fission energy safer?
10.15 Terrain Dependent Renewable Energy Sources
Many renewable energy sources may be installed only in specific areas or terrain;
for example wind energy may be trapped only in those regions where winds of some
minimum velocity below at all times. Same is true for geothermal energy, energy
from ocean and hydroelectric energy. This is partially true for solar energy also;
energy from sunshine may be trapped and converted into electric energy only where
it is going waste, like on rooftops, deserts, etc. Sunlight falling on agricultural land,
in forest areas, etc. cannot be obstructed as it is already being used in such areas by
crops and plants/trees.
10.15 Terrain Dependent Renewable Energy Sources 515
10.15.1 Geothermal Energy
It is simply the power derived from the earth’s internal heat. Geothermal energy is
contained in the rocks and fluids beneath the crust of the earth. At some locations
it may be available at shallow grounds, while at other places it might be few kilo-
metre beneath the surface. Geothermal energy finds its way to the earth’s surface in
three ways: hot springs, geysers and fumaroles (where volcanic gases are released),
volcanoes. Most of the active geothermal reserves are found along major tectonic
plate boundaries, the most active geothermal area in the world called the ring of fire,
encircles the pacific ocean. When magma reaches near the surface of the earth it heats
groundwater trapped in porous rocks or water trapped in fractured rock surfaces and
faults making geothermal water reservoirs.
In USA most of the geothermal power plants are located in western states and
Hawaii where geothermal resources, mostly geysers, are near to the ground level.
California State of USA is producing electricity from geothermal resources since
1960. The first geothermal electric power station was built in Larderello, Italy, in
1904. Geothermal power plants may be of three types: dry steam, flash and binary.
In dry steam plants steam is taken out of the faults or fractures on the ground and is
directly used to run electric turbines. Flash plants pull out hot and high-pressure water
into cooler and low-pressure water producing steam which is used to run turbines.
In binary plants hot water pulled from geothermal reservoir is brought in contact
with some fluid that has low boiling point. The low BP fluid gets vaporised, and
these vapours are used to derive turbines. Many countries including New Zealand,
Iceland, Philippines, Indonesia, Mexico, Sweden and Turkey are producing electric
power using geothermal energy sources.
Geothermal energy is a relatively clean source, energy can be extracted without
burning fossil fuels, and a geothermal power plant emits almost one sixth amount
of CO2 compared to a clean natural gas fired plant of same capacity. As a matter
of fact binary geothermal plants do not release any carbon dioxide. In many cases
geothermal energy may be directly used for heating of homes, removing snow and
in many similar other applications. Geothermal sources supply energy continuously
are not weather or time dependents which is not the case with wind or solar energy.
Figure 10.21 shows a geothermal electric power station.
Many countries have considerable geothermal potential; however electricity
production from geothermal sources is not economically viable at present.
The main concern of geothermal energy is the hydrogen sulphide (H2 S) gas which
is invariably released from geothermal sources. At many geothermal sources a toxic
fluid is also emitted disposal of which creates problems.
Fig. 10.21 Typical

geothermal power station in
New Zealand
10.15.2 Hydroelectric Energy
Hydroenergy, energy obtained by converting the kinetic energy of moving water

mass, is one of the oldest and widely used sources of energy. In almost all earlier
civilisations water flowing in rivers, dropping down in the form of waterfalls, etc.,
has been used to run mills for crushing grains and for other purposes. In 1880 the
first turbine driven by water was used to generate electricity to light an arc lamp,
and in 1881 a dynamo connected to a water-run flour mill provided street lighting
at Niagara Falls, USA. Since then more than 160 countries in the world are now
using hydroelectric power at some level in their basket of energies. World’s biggest
hydroelectric plant is the Three Gorges Dam Hydroelectric plant built on River
Yangtze in China with an installed capacity of 22,500 MW. Figure 10.22 shows a
typical hydroelectric power station.
Hydroelectric energy is one of the most reliable, renewable and low greenhouse
gas emitting source, but it has its own pitfalls, and some advantages and disadvantages
of hydroelectric power are listed here.
Fig. 10.22 Hydroelectric

power station
10.15 Terrain Dependent Renewable Energy Sources 517
View of a typical hydroelectric power station is shown in Fig. 10.22

(a) Advantages
(i) It is a time tested, reliable, adjustable source of base load electrical power
that does not depend on the day or month of the year. In most hydroelectric
power stations, river water is accumulated in a reservoir often made by
building dams. Dams are provided with gates through which accumulated
water is made to fall and run the turbines. Controlling the number of open
gates, it is possible to control the production of electricity. Control on
power generation is a unique feature of hydroelectric power station which
other power generating systems do not have. Electric power generation
may be increased during the time when demand is at its peak and may be
reduced when demand is normal. Since power generation is continuous,
hydroelectric power may be used to run industries without any break.
(ii) Dams built for hydroelectric projects also supply water for irrigation
during the periods when other natural water resources go lean.
(iii) Dam-building activities help in developing the infrastructure in the area.
(iv) It is a renewable source, and rivers which feed reservoirs flow round the
year and are expected to continue like that for a long time to come.
(v) Process of electricity production in hydroelectric power project does not
emit any greenhouse gases and therefore it is a reliable and emission free
source of energy. However, greenhouse gas emission does take place in
construction of dams and other buildings.
(vi) Hydroelectric power in long run is one of the cheapest sources of electrical
energy. High initial cost of dam construction, etc., gets even out in long
run. Normal life of a hydroelectric project is around 30–50 years, but they
continue without much repairs, etc., for longer periods of time.
(b) Disadvantages
(i) Hydroelectric power can be generated only at some places, where there
is a river that flows round the year and has suitable land around available
for building dam and other similar structures.
(ii) Generally large land area, sometimes human and animal habitats, got
submerged in the reservoir, displacing large human population and
destroying fauna and flora of the region.
(iii) Greenhouse gases are emitted in considerable amount in the initial stages
when dam and other power house buildings are constructed.
(iv) Once the flow of a river is restricted, riverside habitats begin to disappear,
adversely affecting the animals and other life forms.
(v) Accumulation of huge quantity of water in reservoirs makes flash flood
threat a reality. Largest number of human and animal lives has been lost
in floods caused by the collapse of dams world over.
10.16 Wind Energy
Kinetic energy of naturally blowing wind may be trapped using a turbine that converts
it into electricity. Special turbines used in wind-run power generation may be clas-
sified into two types: the horizontal-axis wind turbine (HAWT) and the vertical-
axis wind turbine (VAWT). Blade design of the two types of turbines is shown in
Fig. 10.23.
A single wind turbine may provide electric power to a single house, but a cluster
of wind turbines, called a wind farm, may provide electric power to a city. One big
problem with wind power is that it depends on the availability in a given area of
wind with a minimum speed of 12–14 km/h and up to the speed of 50–60 km/h when
electricity generation is a maximum. Further, at wind speeds of around 90 km/h the
turbines must be switched off to avoid damage. The upper and lower limits of wind
speeds essentially depend on the type and size of the turbine blades. Since the output
power level of a wind turbine depends on the speed of the turbine rotor which varies
with the speed of the wind, most of the wind turbines are attached to storing batteries
to stabilise the inconsistent energy surges to be useful. Storing batteries attached to
a wind turbine also provide power at a constant level when the wind speed at wind
farm area is either less than the minimum or is above the maximum speed limit
recommended for the system.
Wind energy is renewable, free from any greenhouse emission and is also cheap.
Wind power share of worldwide electricity usage in 2014 was around 3.1% which
rose to 4.8% by the end of 2018. Portugal and Spain both produce around 20% of
their total energy through wind. India is perhaps the only developing country where
wind energy share is around 10% of total energy consumption, followed by USA
with 10.2% and China 6.1%.
Average cost of wind energy in USA varies between 1 and 2 cent per kWH. On the
other hand in India it is around INR 2.77 (3.5 US cent), in Italy around 7.5 US cent,
Fig. 10.23 Design of the

blades of the horizontal-axis
and vertical-axis wind
turbines and the block
structure of horizontal axis
wind turbine
10.17 Solar Energy 519
and Europe average is around 6.2 US cent per kWH. The global weighted average
of electricity of new onshore wind farms in 2019 was USD 0.053/kWh.
The biggest limitation of wind power is that it requires large open area where
natural wind with speeds above the minimum requirement is available for consider-
ably long periods of time. Another fact is that generally areas where wind power may
be harnessed are far away from cities and places where energy is used. Transportation
of energy adds to the cost of wind energy. Large numbers of wind turbines in wind
farm adversely affect the birds and other wild life in the area.
SAQ: A wind field power station is to be connected to an AC domestic grid. Name
the electronic circuits/blocks that may be required for this.
10.17 Solar Energy
Energy coming in the form of light and heat from sun is called solar energy. Solar
energy may be directly used for heating and cooking in solar heaters, or it may be
converted into electricity. Solar energy is also used to create renewable fuels like
hydrogen. By the end of 2020, around 3% of global energy demand (≈ 700 GW)
was met by solar energy. Further, solar energy sector is the fast-growing sector of
renewable energy. One reason for this rapid growth in solar electric energy use is the
sharp drop of the global levelised cost (LCOE) of solar electricity which dropped by
a factor of 85% during the period from 2010 to 2011.
Solar energy may be used through two different technologies: solar thermal and
solar photovoltaic (PV).
10.17.1 Solar Thermal
In solar thermal technology, radiations from the sun are converted into heat energy
which may be directly used for heating, or it may be converted to steam by concen-
trating incident solar radiations. The steam may then be used to run electric turbines
for generating electricity.
The small-scale thermal technology is used to heat a volume of space, like
rooms and houses and for providing hot water for homes and swimming pools, etc.
The concentrated solar thermal technology is used to produce electrical energy by
concentrating solar radiations falling on a large area. The average power density of
sun radiations on the surface of earth is approximately 1.4 kW per m2 and is rather
small. Solar power may be concentrated at a small area or volume by accumulating
solar powers falling at different areas using reflecting mirrors. The small area or
volume where solar energy is concentrated by reflections from mirrors is called the
receiver. Receiver is kept in thermal contact with a heat reservoir or thermal energy
storage system. Heat energy may be drawn from the heat storage system as and when
required.
As already mentioned the energy density of solar radiations is rather small, there-
fore, it is required to collect solar radiations falling on different areas of the surface
to concentrate large amount of solar energy. This is called solar energy concentra-
tion. In older version of solar energy concentrators, called parabolic concentrator
(see Fig. 10.24), a large parabolic shaped reflecting surface was used. Solar radia-
tions falling on different parts of the parabolic surface got reflected by the surface
to the focus of the parabola and thus large amount of solar energy got deposited at
the focus of the parabola where some energy absorber like water (called receiver)
etc. may be kept. In parabolic concentrators temperatures as high as 400 °C have
been produced at the receiver of the device. Since it is difficult to make very large
parabolic reflecting surface, parabolic concentrators are limited in their applications
and have now been replaced by another kind of concentrator called Tower Concen-
trator. A tower concentrator uses large number of mirrors inclined in such a way
that reflected solar radiations from all mirrors get collected in a small volume at a
tower which works as the receiver. Working principle of a tower concentrator is also
shown in Fig. 10.24. Since there is no limit on the number of mirrors that may be
used to concentrate solar radiations, the total solar energy absorbed by the receiver
may be very large, temperatures of the order of 565 °C at the tower concentrator
receiver has been observed and solar radiations falling on areas of the order of 1–2
Km2 have been concentrated at the receiver. The receiver may be attached to a water
line that may work as the heat store converting liquid into vapours and running the
electric turbine. It is also possible to attach receiver with molten salt storage system
allowing the system to operate for periods of low or no solar energy. Solar power
tower at Crescent Dunes solar energy project concentrates solar energy using 10,000
mirrored heliostats spread on an area of about 1.21 km2 . Parabolic concentrators,
used earlier collect solar energy from a relatively small area.
Fig. 10.24 Tower and parabolic concentrators

10.17 Solar Energy 521
Concentrated solar power (CSP) had a global capacity of around 7000 MW,
the maximum share of about 33% being produced in Spain. Other CSP generating
countries include USA, North Africa, India and China.
SAQ: What is the advantage of tower concentrator over the parabolic one?
10.17.2 Solar Photovoltaic (PV) Technology
Photovoltaic devices convert solar energy into electrical energy. A single PV device,
called a cell, is small, usually less than the thickness of four human hairs, and may
typically produce about 1 or 2 watts of electric power. PV cells are delicate, and
in order to withstand the outdoors for many years, they are sandwiched between
protective transparent materials like in a combination of glass and plastic. To boost
the power output of PV cells they are connected together in chains to form larger
units called panels or modules. Individual panels may be directly used, or they may
be connected to form arrays. Two other important components of a PV power system
are: (i) the mechanical mounting on which arrays are placed to face the sunlight. (ii)
The electrical out of PV array is in the form of DC power, and it must be converted
into AC power before it is connected to the power grid. Alternator circuit converts
DC power into AC power.
Figure 10.25 shows the structure of a PV cell. As shown in the figure a PV cell
is fabricated by developing a very thin layer of N-type material followed by a thick
P-type layer on the same semiconductor crystal. The semiconductor crystal behaves
as an n-p junction diode, with a depletion layer sandwiched between the N- and the
P-sides. It may be recalled (see Chap. 3) that the depletion layer of an np junction
behaves as a charged parallel plate capacitor with N-layer behaving as positively
charged plate and the P-layer as negatively charged plate of the capacitor. Further,
the depletion region is free of any free charge carriers: electrons and holes. It may
also be recalled that there is an internal potential difference V ib between the N-
and the P-sides, N-side being at a higher potential than the P-side. The N-layer is
made thin so that sunlight photons falling on it pass through it. Some conducting
strips are also developed on the N-layer, which work as electrode for taking electric
connection. The P-layer is sufficiently thick so that sunlight photons passing through
the thin N-layer are all stopped in the depletion region and may not leave the cell. An
electrode, for electric connections, is developed on the thick P-layer. The top side
of the N-layer which is exposed to the sun is painted by an antireflection paint so
that most of the sunlight photons falling on the cell are transmitted to the N-layer
without much reflection. When PV cell is exposed to sunlight, N-side facing the
sun, sunlight photons reach the depletion layer passing through the N-layer. Sunlight
photons ionise the atoms of depletion layer creating free charge carrier: electron and
hole. Electrons are attracted by the N-layer which is at positive potential (+V ib ) and
behave as a positively charged plate, while positively charged holes move towards
the negatively charged P-layer at lower potential. The flow of electrons and holes
(created by sunlight photons) constitutes a current through the external load in the
Fig. 10.25 Structure of a

photovoltaic cell
circuit. In this way the solar energy which is contained in sunlight in the form of
photons gets converted in to electrical energy by the PV cell.
Most of PV cells are Silicon cells, which have different conversion efficiencies and
costs ranging from amorphous Silicon cells to polycrystalline and monocrystalline
Silicon types. The efficiency of an ordinary PV cell ranges from 10 to 20%, because of
several factors including the energy of sunlight photons, their absorption in depletion
layer, cleaning of panel surface, etc. Solar PV cell arrays may power a small laptop
to a house. Solar farms of say of 20,000 panels spread across an area of 30 acres
may generate around 3–4.5 megawatts of power, sufficient to power 1200 homes.
Governments in many countries are encouraging house owners to have solar energy
systems for their household use. Household solar systems are very reliable, have low
maintenance cost and have no negative environment impact; if in peak hours the solar
system provides excess electrical energy, and it may be given to be sold to the grid
for use by others. Solar panels are generally put on rooftops, an area that may not
be used for any other purpose. Since most solar cells are made from Silicon, which
is quite in abundance in soil, there is no constrain on the fabrication of PV devices
from the point of the availability of raw material.
A major concern is the high initial cost of fabricating PV cells, panels and arrays.
Fabrication of solar energy systems and solar energy industry adds greenhouse gases,
though use of solar energy in itself is non-polluting and free of greenhouse gas
emission. In view of low conversion efficiency of solar devices coupled with the low
energy density of solar radiations, it is required to spread solar panels or other solar
energy pick-up systems over wide areas to make a viable alternate energy system
based on solar energy. Globally, solar PV electrical generation is expected to grow
by 18% or more in the coming years.
SAQ: Why PV technology for harnessing solar energy is becoming popular?
10.18 Energy from Ocean 523
10.18 Energy from Ocean
Ocean energy may be classified in two broad types: mechanical energy that is asso-
ciated with tides and waves in oceans and thermal energy that is stored by the sun in
ocean in the form of temperature difference between the water at the ocean surface and
deep inside. There is another form of ocean energy the potential energy which is asso-
ciated with ocean currents and is generally not taken into account while discussing
extraction of ocean energy as a source of alternate green energy.
10.18.1 Tidal Energy
Tides in oceans are produced by the gravitational pull of sun—moon–Earth system

on the water of oceans on earth. Tides are long-period waves that result in cyclic rise
and fall of ocean’s surface together with horizontal currents. Tidal energy may be
harnessed using three different technologies: streams, barrages and lagoons.
A tidal stream is fast-flowing body of water caused by tide. Electrical turbines
are directly placed in tidal streams to convert kinetic energy of stream into electrical
energy. The energy density in streams is much larger than that of wind or solar energy
density, and therefore the power output of stream turbines is quite appreciable. Since
duration of tide and the parameters associated with streams at a given place are all
quite predictable, electrical energy produce through tidal streams is quite steady and
reliable as compared to the wind energy. However, there are few problems associated
with the placing of turbines in tidal streams; firstly, placement of turbines interfere
with the tide and may reduce the intensity or alter the direction of tidal stream;
secondly, if turbines are placed in the deeper part of the stream, they may not be
visible and navigation around them may be disrupted. And thirdly, marine life may
be severely endangered by the rotating blades of the turbine. It has been observed that
best results are obtained when turbines are placed in shallow part of the stream since
turbines are visible and navigation around them is possible. Moreover, in shallow
part of the stream velocity of stream flow is less, and hence the turbine blades rotate
slowly reducing the danger level to marine life. World’s first tide stream-based electric
power station was built in Northern Ireland at Strangford Lough in 2007 where the
stream velocity was about 4 m/s.
Barrage is a dam of considerable length but of small height that encloses a large
tidal basin and made in the tidal zone to harness tide energy. In high tide, tide water
spills over the barrage or is transported behind the barrage through sluice gates in the
barrage. Water accumulated behind barrage is released through low-level turbines
that convert potential energy of accumulated tide water into electricity. Barrage may
be built across tidal rivers, bays and estuaries.
Barrage system has negative impact on environment. Free movement of fishes
from ocean to barrage basin is restricted, changing water level in basin area disrupt
fauna and flora of the area, and since large part of the basin remains under water it
does not produce food for birds of the area and they migrate. Level of salinity of
the area is also affected. Rotating blades of turbines also endanger marine fife. A
barrage-type power station built at Rance River estuary in Brittany, France, in 1966
to harness tidal energy is still working. It uses the energy contained in currents of the
river as well as the tidal energy of English Channel. The barrage has increased the
silt level and has adversely affected the plant life of the area. A fish of the area called
Plaice has totally extinct, while another variety called cuttlefish is now thriving in
the area.
Tidal lagoon is a body of ocean water that is partially enclosed by artificial or
natural boundary. Tidal lagoon works in the same way as a barrage basin. Ocean water
enters lagoon during high tide and drains out during the low tide. Electrical turbines
generate electricity during both cycles, and therefore, lagoon-based generators may
supply continuous electrical energy.
World’s five biggest power stations based on tidal power are: (i) Sihwa Lake
Tidal Power Station, South Korea, capacity 254 MW; (ii) Rance Tidal Power Station,
France, capacity 240 MW; (iii) MeyGen Power Station, UK, capacity 6.0 MW; (iv)
Jiangxia Tidal Power Station, China, capacity 3.2 MW; and (v) Eastern Scheldt
Barrier Tidal Power Station, Netherlands, capacity 1.25 MW.
Ocean energy is renewable and clean and qualifies to be a sustainable source of
energy; however, the market for tidal electrical energy is not yet grown, and investors
are not sure about its profitability.
10.18.2 Ocean Thermal Energy
Thermal energy in the form of temperature gradient with depth in ocean water is
deposited by the sun. Since the ocean water temperature gradient is in general small
and varies with the location of the place, the first task is to find a suitable site for
installing a system that may harness ocean thermal energy. Ocean thermal energy
conversion (OTEC) plants can only be installed at seashores where the top ocean
surface remains hot all through the year and the water at depth remains sufficiently
cold. Hawaii and some Caribbean islands are most suited for OTEC. However,
these areas are also very rich in solar and wind energies, and therefore the ocean
thermal energy has to compete with these other sustainable energy sources to become
economically viable.
Figure 10.26 shows an artist’s view of an OTEC system. Ocean surface water
at temperature 20–25 °C may be used to vaporise another low boiling point fluid,
called working fluid, and pressurised working fluid vapours may run an electric
generator. Low-pressure and low-temperature vapours discharged from the turbine
(as it converts thermal energy of vapours into electrical energy) may be condensed
using cold water (3–7 °C) pumped from the depth of the ocean and may be returned
to the boiling pot forming working fluid closed loop. Sea water itself may be used as
the working fluid and may be turned into vapours by supplying additional heat. The
advantage of using sea water as working fluid is that the water left as discharge from
10.19 Portable Sources of Sustainable Energy 525
Fig. 10.26 Artist’s design of OTEC
the electric turbine is pure salt-free water and may be used for human consumption
or may be further cooled using the low-temperature ocean water pumped from the
depth, to supply as refrigerated cold water.
The global potential of OTEC is much larger than other ocean energy options,
and it may become a major source of continuous, clean, renewable energy that may
sustain base load and therefore many countries including Japan, USA, European
countries, India, China, etc. are working on projects involving OTEC.
SAQ: What are the three different technologies used to harness tidal energy?
10.19 Portable Sources of Sustainable Energy
Intermittent renewable sources, like the wind electric turbines or solar PV cells, etc.,
which generate electric power only at some particular time or generate excess power
at some peak hour, need a system where excess electrical energy generated may be
stored for use at a later time. Chargeable lithium-ion batteries and supercapacitors are
devices that may be employed to store electrical energy for use at some other place
and time. Both of them are portable and can be carried from one place to another;
not only that, these devices may be refilled with electrical energy once the energy
already stored is consumed.
10.19.1 Lithium-Ion Battery
Lead-acid, Nickel-metal hydride, Nickel-Cadmium and lithium-ion batteries are all

rechargeable batteries in which some reaction converts chemical energy into elec-
trical energy and the reverse reaction converts electrical energy back into chemical
form. Once charged by electrical energy, these batteries may hold the electrical
energy in chemical form and will supply the stored energy back into electrical form
when required. While considering the merits of a rechargeable battery that may be
transported from one place to another, following characteristics become important:
(i) Gravimetric Energy Density It is the ratio of electric energy contained in
the battery to the mass of the battery. A battery with higher energy density is
light in weight with higher electrical energy, may be easily transported and
may be kept in a moving vehicle to provide it power without much increase in
the weight. Typical values of energy density of Li-ion batteries range 100–265
Wh/kg. This value is much higher to the energy density of other rechargeable
batteries like Ni-Cd battery energy density for which is 45–80 Wh/kg.
(ii) Maintenance Li-ion batteries do not require any recycling or periodic topping-
up of acid, etc. which is necessary for the maintenance of other similar type of
batteries.
(iii) Longevity Longevity of a battery is measured in two ways: (a) how many cycles
of charging and discharging the battery may sustain; for example, a lead-acid
battery may withstand 500–1000 cycles, while a Li-ion battery may withstand
up 10,000 cycles. (b) Self-discharge rate tells how a fully charged battery loses
its charge when not in use. Li-ion battery has very low self-discharge rate of the
order of 1.5–2% per month as compared to 10–15% per month for lead-acid
battery. It may however be remarked that self-discharge rate depends on the
environment and is generally low if the ambient temperature is low.
(iv) Voltage Per Cell Batteries are made by joining several basic units called cells.
Each cell of a given battery develops a fixed voltage; for example each cell of
Ni–Cd cell may have up to 2.0 V. Each cell of Li-ion battery develops 3.6–4.2 V
of potential. A higher potential per cell reduces the number of cells in a given
battery pack to achieve the given output. A single Li-ion cell may run a watch
for a year or more.
(v) Charging Speed and Ease of Charging Typical charge and use cycle of a
lead-acid battery is 8 h of use, 8 h of charging, 15 min of cooling as charging
produces substantial amount of heat. On the other hand, the cycle for Li-ion
battery is: 8 h of use, 1 h of charging, 8 h of use; no cooling is required as
negligible heat is produced during charging. However, heat production both
during charging and during use depends on the design of the cell and available
heat dissipation.
(vi) Stability and Safety One limitation of Li-ion battery is that strict control on
maximum value of output current and temperature has to be maintained for
steady and stable operation of the battery. The level of charge to which the
battery may be charged and the level to which the battery may be allowed
to be discharged are all very critical. For long life and safety, it is required
that battery should not be charged beyond 80% of the maximum value and,
similarly, should not be allowed to discharge below 20% value. These controls
in a Li-ion battery are achieved by the use of IC-based electronic circuits.
Figure 10.27 shows the diagram of a Li-ion cell. The cell may be divided into
five components: Aluminium current collector cathode, lithium metal oxide, porous
separator that allows the migration only of lithium ions and does not allow migration
of electrons through it, lithium-loaded graphite and a copper anode. The cell is filled
with a polymer gel/organic carbonate, like ethylene carbonate which serves as the
electrolyte for the migration of lithium ions. When a charged Li-ion cell is connected
to the load, Li-ions from the anode side migrate through the separator to the Cathode
side inside the cell, while electrons travel through the external load constituting the
load current. Reverse happens during the charging phase, and Li-ions migrate from
cathode side to the anode side via separator, and electrons from the charging current
source reach anode to complete the circuit. As already mentioned, the electrolyte
simply facilitates the migration of lithium ions from one side to the other across the
separator.
Fig. 10.27 Layout of a lithium-ion cell

10.19.2 Super Capacitor
Super- or ultracapacitors are used for energy storage undergoing frequent charge and
discharge cycles at high current and short duration.
Supercapacitors are used in automobiles, buses, trains, cranes, elevators and other
airport equipment where regenerative breaking, short-term energy storage or burst-
mode power delivery is required. In contrast, Li-ion batteries are used where long-
term energy storage is required. The layout of a typical Helmholtz double-layer
supercapacitor that develops two layers of positive ions and electrons at the interface
of electrolyte with anode and cathode is shown in Fig. 10.28. As shown in the figure,
a perforated separator layer that allows the migration only of positive ions divides
the capacitor into two sections. Each Helmholtz layer may be treated as a parallel
place capacitor with large quantity of charge on plates and extremely small distance
between the plates. Therefore the capacitance of each layer is quite large in the range
of Farad. The energy storage mechanism in a supercapacitor is bulk separation and
movement of charges.
Supercapacitors have two electrodes: a separator and an electrolyte. Outer sides
of electrodes are covered with Aluminium foils that work as current collectors.
Electrodes are porous generally made of different types of carbon such as carbon
cloth, activated carbon or carbon nanotubes or of mixed metal oxides or conducting
Fig. 10.28 Layout of a double-charge layer supercapacitor

polymers. Recently graphene because of its high chemical stability, high electrical
conductivity and large surface area has also been used for making electrodes. The
separator allows the migration of ions through it but restricts the migration of elec-
trons. The electrolyte is mostly liquid through which ions may easily flow between
electrodes. Energy density of a supercapacitor is quite high, comparable to a Li-ion
cell.
SAQ: Compare a supercapacitor with a lithium-ion cell. The working voltage of a
1F supercapacitor is 3 V; what maximum charge it may hold?
SAQ: What are some special features of Li-ion batteries? List two drawbacks of
Li-ion batteries.
SA10.1 What is meant by sustainability? Discuss how an individual may
contribute to the sustainability of earth.
SA10.2 Discuss the concept of sustainable economy.
SA10.3 What is greenhouse effect and what is responsible for it in earth’s
atmosphere?
SA10.4 What are the five dimensions of social sustainability?
SA10.5 What are the main requirements for economical sustainability? Comment
on circular economy.
SA10.6 What is greenhouse effect? What causes it?
SA10.7 What is the big achievement of Paris Agreement? What is ‘Nationally
Determined Contributions (NDCs) in this context?
SA10.8 What may be the possible ways of removing CO2 from the atmosphere?
SA10.9 Why CO2 is a greenhouse gas and why is it most fatal from the point of
global warming?
SA10.10 What may be the possible ways of removing CO2 from the atmosphere?
SA10.11 Define the following terms and discuss their significance; Primary energy,
LCOE, Gravimetric energy density.
SA10.12 What is a fuel cell? How does it differ from a battery?
SA10.13 What are the main components of a fuel cell? Explain the working of a
(PEM) cell.
SA10.14 How does an (AFC) cell differ from a (PEM) cell? List one advantage and
one drawback of (AFC) cell over the (PEM) cell.
SA10.15 What fuel is used in a (PAFC) cell? Why this cell is often used in big
structures like trucks, buildings, etc.?
SA10.16 Compare advantages and disadvantages of hydroelectric power.
SA10.17 Compare wind, solar and geothermal electric sources of energy
SA10.18 Name different techniques used in trapping ocean energy, explaining one
of them in some details.
SA10.19 With the help of a diagram give the construction details of a lithium-ion
cell and list its advantages.
SA10.20 What is the purpose of a supercapacitor? Explain the formation of

Helmholtz double layer in supercapacitor.

native; all correct alternatives must be marked for complete answer of the question
in such cases.
MC10.1 Steam methane reforming is the commercial method of producing
(a) Carbon dioxide (b) nitrogen (c) oxygen (d) hydrogen
ANS: (d)
MC10.2 The fuel that has very high gravimetric energy density and smallest
volumetric energy density is
(a) Hydrogen (b) nuclear fuel (c) coal (d) petrol
ANS: (a)
MC10.3 Which of the following has the highest energy density?
(a) Hydrogen fuel cell (b) li-ion cell (c) supercapacitor (d) lead-acid cell
ANS: (a)
MC10.4 Which of the following has the highest power density?
(a) Hydrogen fuel cell (b) li-ion cell, (c) supercapacitor (d) lead-acid cell
ANS: (c)
MC10.5 Average energy releases per fission of 235 U nucleus is approximately
(a) 50 meV (b) 200 meV (c) 3 × 10−11 J (d) 10.0 × 10−18 kilowatt hour
ANS: (b), (c), (d)
MC10.6 Which precious metal is used as a catalyst in fuel cells?
(a) Gold (b) Silver (c) copper (d) platinum
ANS: (d)
MC10.7 Sustainable energy source must be
(a) Renewable (b) should not emit any greenhouse gas (c) should be
available for long time to come (d) may emit greenhouse gases in limited
small amount
ANS: (c), (d)
MC10.8 Which of the following energy source has minimum value of LCOE?
(a) Nuclear (b) hydrogen fuel cell (c) wind (d) geothermal
ANS: (c)
MC10.9 Which of the following energy sources may provide electricity and cold
pure water as a by-product?
(a) Hydrogen fuel cell (b) ocean thermal (c) nuclear (d) geothermal
ANS: (b), (d)
MC10.10 Which type of economy is best suited for sustainable growth?
(a) Circular economy (b) free market economy (c) command economy
(d) mixed economy
ANS: (a)
LA10.1 Discuss the concept of sustainability and its components. Give a detailed
account of economical sustainability.
LA10.2 What is global warming? Discuss its causes. Mention about important steps
taken by the UNO in arresting temperature rise.
LA10.3 What should be the essential properties of a sustainable energy source?
Give a list of energy sources that in your opinion are best for your country,
and justify your choice.
LA10.4 Differentiate between a battery and a fuel cell. List different types of fuel
cells, and explain their relative merits and de-merits.
LA10.5 Write a detailed note on solar energy as a viable and sustainable source.
LA10.6 Give detailed account of the construction and working of a lithium-ion cell.
What are the main advantages of this cell?
LA10.7 Write an essay on sustainable energy options.
Index
A Characteristic X-rays, 193, 198–201, 203,

Absorption force, 450 204, 264
Acceptance angle, 358–362, 374, 375 Circular economy, 476, 529, 531
Acceptor, 81, 89, 90, 103–105, 108, 109, Cladding, 349–351, 353–355, 357–362, 428
111, 131, 134 Classical physics, 30, 149, 191, 210, 222,
Active medium, 387, 389–391, 393–396, 223, 227, 228, 240, 248, 249, 265,
401, 403, 405, 406, 411, 415–417, 267, 268, 295, 300, 306, 308, 310,
422–424, 428, 431 317
Alcoysis, 455 Coercive field, 174, 175, 183
Anti-commutator operator, 278 Compensated semiconductor, 90
Aspect ratio, 467, 468, 470 Compressive strength, 10, 25, 28
Assemblies, 317, 437, 457, 470 Compton equation, 244
Atomic specific heat, 247–253, 256 Conductance, 11, 62
Auger effect, 235 Conduction band, 64–69, 73–75, 78, 85–88,
90, 91, 96–98, 102, 108–110, 114,
115, 162, 412, 441, 442
Cooper electron pairs, 127, 128
B
Criticality parameter, 509
Bandwidth, 348, 369, 373, 375, 376
Cross-talking, 348
Barrage, 523, 524
Crystal lattice, 8, 9, 13, 14, 16, 75, 77, 82,
Barrier tunnelling, 300, 301
83, 85–90, 92–95, 113, 126, 127,
Biot-Savart law, 139, 187
172, 206, 207, 253, 443
Blackbody, 191, 235–240, 242, 253, 267, Crystallinity, 1
384 Curie temperature, 164, 165, 168, 170, 172,
Bottom-up, 435, 436, 453, 454, 458, 470, 173, 178, 183, 189
472 Cut-off frequency, 226
Bound states, 129, 290, 291, 295, 313
Breakdown voltage, 67, 107
Bremsstrahlung X-rays, 192, 193 D
Buffer jacket, 349 Degenerate, 90, 277, 283, 341, 345
Degenerate semiconductors, 90
Delocalised electrons, 6, 8–11, 13, 15, 52,
C 53, 59, 75, 78, 441, 443
Calcined, 175, 455 Depletion layer, 103–107, 131, 133, 134,
Canonical conjugates, 308 371, 412–414, 521, 522
Carbon footprints, 485, 486, 495, 511 Di-atomic molecule, 8, 64, 65
© The Editor(s) (if applicable) and The Author(s), under exclusive license 533
to Springer Nature Switzerland AG 2023
https://doi.org/10.1007/978-3-031-32084-2
534 Index
Diffusion coefficient, 84, 95 Greenhouse effect, 477, 478, 480, 482, 484,
Diffusion currents, 94 529
Direct semiconductors, 91, 412 Greenhouse gasses, 478–488, 492–494,
Distribution function, 97, 240, 242, 272, 501, 503, 516, 517, 522, 529, 530
318, 338, 339, 342, 344, 346 Group velocity, 217, 219–221, 293, 294,
Donor, 81, 85–87, 90, 103–105, 108, 109, 365, 367
134 Gyromagnetic ratio, 145, 149
Doped semiconductors, 68, 81, 90, 108, 110
Drift velocity, 93, 100, 113, 443
H
Half metal, 68, 69, 114, 115
E Hall effect, 99, 100, 131, 135
Economic sustainability, 475, 476 Hermitian operators, 267, 268, 270, 275,
Eigenvalue, 271, 276, 277, 283–287, 297, 276, 281–283, 286, 312–315
312 High dimension arrays, 437
Einstein’s specific heat equation, 251 Hole, 74, 75, 78, 80, 81, 86–89, 91, 92,
Electric dipole, 51, 53, 54, 141, 142, 451 94–96, 100, 102–105, 107–111, 114,
Electronegativity, 44–46, 52, 55, 56, 59, 60 132, 236, 372, 405, 412, 414, 441,
Electron mobility, 93 442, 521
Energy carrier, 502 Homogeneity, 272
Ensemble, 317, 437
Environment, 71, 74, 181, 373, 458, 461,
I
473, 474, 476, 477, 481, 484–487,
Impact strength, 10, 24
489, 493, 495, 501, 522, 523, 526
Indirect semiconductors, 90, 91
Environmental sustainability, 476, 477
Induced decay, 384
Exchange interactions, 165, 168, 170–172,
Insulators, 2, 28, 61, 63, 64, 66–68, 75, 78,
176–178, 186, 188, 189
85, 104, 131, 134, 138, 175, 182,
Expectation value, 286, 287, 289, 310, 313,
263, 406, 408, 440, 443
315
Intermetallides, 19
Internal potential barrier, 103–107
Intrinsic semiconductors, 68, 73–75,
F
78–81, 85–90, 101, 108–110
Fermi Sea, 72
Ferrites, 18, 19, 181–184, 187, 189
Fick’s second law, 84 L
Fluorescence, 415, 416, 438, 449 Langevin function, 160, 161
Food deserts, 474 Leveilized cost of electricity, 499
Forbidden energy gap, 64–69, 73–75, 78, Lewis symbols, 42
86, 89, 91, 109, 110, 113–115, 129, Lithography, 459
130, 440 Localised Surface Plasmon Resonance,
Forward bias, 105, 107, 412, 413 440, 470
Fuel cell, 468, 503–506, 529–531 Lustre, 2, 3, 10
Full angle divergence, 403
Functionality, 21, 60
M
Macromolecules, 19
G Macrostate, 323–328, 331, 333, 336,
Gamma rays, 210, 216, 242, 243, 264, 378 340–346
Garnets, 181, 182 Magnetic spin quantum number, 30, 31, 33
Gelation, 455 Majority carriers, 86–90, 96, 103, 106, 107,
Giant pulse formation, 423 112, 133, 412
Gradient force, 450–453, 470, 472 Malleability, 3, 5, 9
Gravimetric energy density, 499, 500, 502, Matrix, 24–28, 58, 59, 184, 268, 275, 276,
505, 526, 529, 530 467, 505
Index 535
Meissner effect, 118, 123, 134 Population inversion, 377, 386, 388–394,
Metallic bonding, 6, 7, 9, 18, 52, 113, 115 402, 406–409, 411, 412, 414, 416,
Metalloids, 69, 114 418–422, 431, 433
Metastable states, 389–392, 401, 402, 409, Primary carbon footprint, 485
416, 418, 431 Primary energy, 495, 496, 501, 529
Microstates, 323–327, 334, 336, 340–342, Pumping, 389–392, 394, 395, 401, 402,
344–346 406–411, 415, 416, 420, 421, 424,
Minority carriers, 86–89, 95, 107, 110, 112, 428, 433
133
Mode-locked, 423
Monocrystalline, 69, 71, 522 Q
Quantum confinement, 441, 471
Quantum dots, 437–440, 442, 458, 470, 471
N Quantum mechanical superposition, 272,
Nano films, 437 292
Nanofluids, 444 Quantum mechanical tunnelling, 306, 308,
Nanoscale, 435, 437 315
Nanotechnology, 435, 437, 467, 469, 470
Nano tubes, 437, 439, 464, 466, 469
R
Nano wires, 437, 438, 444, 458
Radiationless de-excitation, 383
Nanoworld, 435, 437
Rayleigh Jean’s formula, 240
Nationally Determined Contributions
Residual magnetization, 174, 175
(NDCs), 487, 529
Resistivity, 11, 14, 19, 56, 61–64, 95, 96,
Non-degenerate level, 320, 330
113, 114, 116–118, 120, 126, 131,
Non-degenerate semiconductors, 90, 108
134, 175, 443, 466
Numerical aperture, 360, 361, 373
Retarding potential, 227, 230
Reverse bias, 105, 107, 134, 372
Right hand rule, 139, 140, 142
O
Occupation number, 319, 321–324, 328,
331, 334, 338, 341–344 S
Optical resonator, 393, 402 Scattering force, 450–452, 472
Schrodinger’s time-independent equation,
274, 289, 313
P Schrodinger wave equation, 30, 268
Pair production, 246 Screening constant, 204
Partition function, 318, 339 Secondary carbon footprint, 485
Pauli’s exclusion principle, 8, 31, 33, 170, Self-assembly, 436, 455, 456, 470, 472
176, 188 Semiconductors, 2, 17, 28, 61, 63, 64, 66,
Permeability, 139, 152, 153, 175, 189 68–70, 72, 74, 75, 78, 79, 81–97,
Phase velocity, 217, 218, 220, 221, 294 99–101, 107–111, 113, 114, 124,
Phonon, 14, 116, 126–129, 135, 172, 253, 130–132, 134, 391, 405, 406, 409,
255, 443 412, 414, 433, 436, 437, 440–443,
Phosphorescence, 416 458, 521
Photo avalanch diode, 372 Shallow impurities, 90, 92
Photon, 8, 11, 14, 52, 91, 110, 130, 191, Shear strength, 10
195, 196, 211, 216, 232–235, Singlet state, 176, 177, 187, 415
242–246, 253, 258–261, 263, 264, Social sustainability, 474, 475, 529
364, 378, 380–383, 385–402, 409, Specific dielectric strength, 67
411–414, 416–418, 423, 429, 431, Stationary states, 274, 310
432, 438, 443, 444, 450, 451, 458, Superconductors, 61, 63, 64, 117–126, 129,
481, 521, 522 131, 135, 157
Planck distribution function, 242 Supermagnetism, 445
536 Index
Surface Plasmon, 439, 440, 470 U

Sustainable development, 473, 488 Ultimate strength, 10
Sustainable energy, 473, 476, 495–497, Ultraviolet catastrophe, 240
524, 525, 530, 531
Sustainable energy sources, 495, 496
Sustainable land management, 489 V
Sustainable society, 473, 474 Valence band, 64–68, 73–75, 87, 89–91, 96,
108–111, 113–115, 412, 413, 441,
442
Vander-Waal bonds, 18, 53
T Volumetric energy density, 499, 502, 530
Tensile strength, 10, 16, 18, 25–27 Vortices pinning, 124
V parameter, The, 362
Thermal cavity, 236
Thermal conductivity, 14, 15, 444, 445,
467, 468, 471
W
Thermal equilibrium, 134, 236, 239, 241, Wavefunction, 247, 269, 271, 272,
379, 384–386, 388, 431, 433, 507 278–280, 283, 291, 292, 294, 298,
Thermal radiations, 235–237, 477, 481 301–303, 308, 310, 312, 313, 315
Tidal lagoon, 524 Wien’s displacement law, 237, 262
Tidal stream, 523 Wien’s distribution law, 237, 238, 240
Top-down, 435, 436, 453, 457, 458, 470, Work function, 227, 230, 232–234, 258,
472 262
Transition temperature, 5, 117, 118, 121,
164, 168
Triplet state, 177, 187 Y
Tuning, 402, 409 Yield strength, 10, 27

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Physics and Technology for Engineers

Physics and Technology

ISBN 978-3-031-32083-5 ISBN 978-3-031-32084-2 (eBook)

(Beautiful lady fought bravely with cancer

Kellyville, NSW, Australia R. Prasad

Kellyville, Australia R. Prasad

1 Engineering Materials, Atomic Structure and Bounding . . . . . . . . . . 1

2.4.5 n and p Type Semiconductors . . . . . . . . . . . . . . . . . . . . 85

4.2.6 Some Application of X-rays . . . . . . . . . . . . . . . . . . . . . . 208

5.5 Time-Independent Schrodinger Equation . . . . . . . . . . . . . . . . . . . 273

6.9 The Distribution Function . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 338

9.4.1 Bottom-Up Techniques . . . . . . . . . . . . . . . . . . . . . . . . . . 453

10.14.2 Plus Points of Fission Reactor . . . . . . . . . . . . . . . . . . . . 511

1.1 Classification of Condensed Matter

© The Author(s), under exclusive license to Springer Nature Switzerland AG 2023 1

Fig. 1.1 Periodic table of elements

are either deliberately mixed in controlled amount or are present as impurity in a

Table 1.1 Density of different materials

is defined as the strain at which a test specimen fractures during a uni-axial

final length − initial length

The percentage reduction in area is given as;

Percentage reduction in area

It is important to note that the property of ductility depends on the tempera-

Fig. 1.2 Ductile and brittle

Fig. 1.4 In case of metal

Fig. 1.6 Closely spaced

Fig. 1.7 A rough and

Strength of a material is judged by the ability of the material to resist deformation

Table 1.2 Tensile and yield

excited delocalised electrons almost immediately revert back to their ground

The constant of proportionality ρ in above equation is called the resistivity

Fig. 1.8 Bar graph showing

Table 1.3 Conductivity for

Fig. 1.9 Current flow

Fig. 1.10 Heat flow across

temperature difference ∆T and is inversely proportional to the thickness ∆x.

A . (∆t) . (∆T ) A.(∆t).(∆T )

Table 1.4 Thermal conductivities of some materials

binding energy of crystal lattices. Tightly bound lattices in metallic crystals

Table 1.5 Melting and boiling points of some metals

It is difficult to define ceramics in the present context as ceramics cover a very

Table 1.6 Broad features of Ceramics

Fig. 1.11 Bar graph

Table 1.7 Effect of bond

Fig. 1.12 a Bond structure

Table 1.8 Variation in some

The difference in mechanical properties arises of their structures. Figure 1.13

Fig. 1.13 Structure of wax

Fig. 1.14 Structure of

(g) Semisynthetic polymers Derivative cellulose that is obtained by modi-

nCH2 = CH2 →→ Polymerization →→ −(CH2 − CH2 )n −

Polymers made by addition polymerisation of two different monomers

Fig. 1.15 An example of

SAQ: Name three types of bonds that are found in polymers.

Fig. 1.16 Four types of loads that may be applied to a composite

Fig. 1.17 Ceramic matrix

1.2 Atomic Structure

1.2.1 Elements of Atomic Structure