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EDUCATIVE COMMENTARY ON

JEE 2023 ADVANCED MATHEMATICS PAPERS


(Last Revised on 15/04/2024)

Contents
Paper 1 2

Paper 2 38

Concluding Remarks 64

This is a commentary on the problems from the Mathematics section of


the JEE 2023 Advanced papers. The 2019 commentary was a departure
from the practice for the then preceding 16 years (from 2003 to 2018) where
the commentaries covered all questions. Instead, a few questions were omit-
ted from the commentary. There were two reasons for this change. First,
some problems are straightforward and similar problems had been asked in
the past. They lacked any significant educative content. Secondly, detailed
solutions to such problems are easily available on many websites and there
is not much point in duplicating them. For example, routine problems in
coordinate geometry and trigonometry were omitted.
Since 2020 this selective practice has been modified. The solutions to all
the problems are included. But when there is not much to comment, only
the answer and the important steps in the solution are given. Hopefully, the
readers can fill the gaps.
I thank those who provided me the question papers. Vardhan Kumar Rai
pointed out the lower bound L(x) on f (x) in the solution to Q.4 in Paper 1.
Abhishek Bansal pointed out that in Q.14 of Paper 1, mere vanishing of the
four determinants does not imply that there are infinitely many solutions.
Readers who notice any errors, want to comment or give alternate solu-
tions are invited to send an email to the author at kdjoshi314@gmail.com or
send an SMS or a WhatsApp message to the author at 9819961036.
Unless otherwise stated, all the references made are to the author’s book
Educative JEE (Mathematics) published by Universities Press, Hyderabad.
PAPER 1

Contents

Section - 1 (One or More Correct Options Type) 2

Section - 2 (Only One Correct Option Type) 9

Section - 3 (Non-negative Integer Answer Type) 17

Section - 4 (Matching Entries in Columns Type) 29

SECTION - 1
This section contains THREE questions each of which has ONE OR
MORE correct options. There are 4 marks if all correct and no other op-
tion(s) are chosen, 0 marks if no answer is given and -2 marks if even one
wrong option is chosen. There is partial credit if some but not all correct
options are chosen and no incorrect option is chosen.

Q.1 Let S = (0, 1) ∪ (1, 2) ∪ (3, 4) and T = {0, 1, 2, 3}. Then which of the
following statements is(are) true?

(A) There are infinitely many functions from S to T


(B) There are infinitely many strictly increasing functions from S to
T
(C) The number of continuous functions from S to T is at most 120
(D) Every continuous function from S to T is differentiable

Answer and Comments: (ACD). The truth of (A) is trivial because


the set S has infinitely many points and we can arbitrarily send each to
any of the four elements of T . Falsity of (B) is only sightly less trivial
because a strictly increasing function is automatically one-to-one and
there cannot be a one-to-one function from an infinite set to a finite one.
(C) is based on the non-trivial Intermediate Value Property. We claim
that a continuous function f on an interval I with only finitely many
values must be a constant. For, otherwise there will be two distinct

2
values, say, α and β which f assumes at two points, say a and b in I
respectively. Without loss of generality assume α < β. As there are
infinitely many points in the interval (α, β), there is some γ ∈ (α, β)
which is not a functional value of f . But that violates the IVP applied
to the interval (a, b) (or the interval (b, a) depending upon whether
a < b or b < a).
It is now easy to show that both (C) and (D) are true. As the codomain
T is finite, a continuous function from S to T must be a constant on
each of the three intervals (0, 1), (1, 2) and (3, 4). These three constants
are independent of each other. As there are 4 possibilities for each, the
number of distinct continuous functions from S to T is 43 = 64 < 120.
The underlying fact about (C) also makes (D) true because on each of
the three intervals, f is constant and hence differentiable.

A simple problem on theoretical calculus. A better question would


have been to ask which property of continuous functions makes (C)
true, the options, in addition to IVP being, say, the existence of a
maximum, or of a minimum and integrability.

Q.2 Let T1 and T2 be two distinct common tangents to the ellipse E :


x2 y 2
+ = 1 and the parabola P : y 2 = 12x. Suppose that the tangent
6 3
T1 touches P and E at the points A1 and A2 respectively and the
tangent T2 touches P and E at the points A4 and A3 respectively.
Then which of the following statements is(are) true?

(A) The area of the quadrilateral A1 A2 A3 A4 is 35 square units


(B) The area of the quadrilateral A1 A2 A3 A4 is 36 square units
(C) The tangents T1 and T2 meet the x-axis at the point (−3, 0)
(D) The tangents T1 and T2 meet the x-axis at the point (−6, 0)

Answer and Comments: A, C.


The data is shown in the figure below. Since both E and P are sym-
metric about the x-axis, it is intuitively obvious that the tangents T1
and T2 are reflections of each other in the x-axis. But we shall not take
this for granted. Rather we shall prove it.

3
y

A.1
T1
P
A2.
E
O x
.
A3
T2 .
A4

The condition for tangency of a line T : y = mx + c for the parabola


y 2 = 12x is
3
c= (1)
m
x2 y 2
while its condition for tangency for the ellipse + 2 = 1 is
6 b
c2 = 6m2 + 3 (2)

For T to be a common tangent, we solve (1) and (2) simultaneously.


Eliminating c between (1) and (2) gives
9
= 6m2 + 3 (3)
m2
which gives a biquadratic in m, viz.

6m4 + 3m2 − 9 = 0 (4)

whose roots are m2 = 1 and m2 = − 32 , As m is real, we have m = ±1


and c = ±3.
So the two common tangents are

T1 : y = x + 3 (5)
and T2 : y = −x − 3 (6)

4
Both T1 and T2 meet the x-axis at (−3, 0). Hence (C) is true.
To find the point of contact A1 of T1 with P , we solve y = x + 3 and
y 2 = 12x. Eliminating x, (x + 3)2 = 12x, i.e. (x − 3)2 = 0. As expected
(because of tangency), there is a double root 3. Then y = 6. Hence

A1 = (3, 6) (7)

By symmetry, A4 = (3, −6). Similar calculations with the ellipse in-


stead of the parabola give A2 = (−2, 1) and A3 = (−2, −1).
The quadrilateral A1 A2 A3 A4 is a trapezium with parallel sides of lengths
2 and 12. The horizontal distance between them is 3−(−2) = 5. There-
fore, the area of A1 A2 A3 A4 is 21 × (12 + 2) × 5 = 35 sq. units. Hence
(A) is also true.

A conceptually simple but computationally long problem based on


condition for tangency.
x3 5
Q.3 Let f : [0, 1] −→ [0, 1] be the function defined by f (x) = −x2 + x+
3 9
17
. Consider the square region S = [0, 1]×[0, 1] . Let G = {(x, y) ∈ S :
36
y > f (x)} be called the green region and R = {(x, y) ∈ S : y < f (x)}
be called the red region. Let Lh = {(x, h) ∈ S : x ∈ [0, 1]} be the
horizontal line drawn at a height h ∈ [0, 1] . Then which of the following
statements is(are) true?

(A) There exists an h ∈ [ 14 , 32 ], such that the area of the green region
above the line Lh equals the area of the green region below the
line Lh
(B) There exists an h ∈ [ 14 , 32 ], such that the area of the red region
above the line Lh equals the area of the red region below the line
Lh
(C) There exists an h ∈ [ 41 , 32 ], such that the area of the green region
above the line Lh equals the area of the red region below the line
Lh
(D) There exists an h ∈ [ 41 , 32 ], such that the area of the red region
above the line Lh equals the area of the green region below the
line Lh

5
Answer and Comments: B, C, D. In the past, because of the lim-
itations of the printing technology, diagrams were generally not drawn
in the question papers, especially in mathematics where the ability to
translate a verbal description to a visual diagram was also a quality
to be tested. Today, things have changed. Diagrams can already be
inserted into the file sent to the printer (or to the candidates in the case
of an online test). And considering the severely limited time allowed for
each question, it is plainly unfair if a candidate has to spend a precious
portion of it just to understand the problem. So a diagram should have
been included as there is plenty of work left even after that.

y
(0,1) (1,1)

2/3

181/324
17/36 green
red Lh
h
13/36
1/4

x
O (1,0)

Let AR and AG denote the areas of the red and the green regions
respectively. Then by a direct calculation,
Z 1
AR = f (x)dx
0
x3
Z 1 5x 17
= − x2 + + dx
0 3 9 36
x4 x3 5x2 17x 1
= ( − + + )
12 3 18 36 0
1 1 5 17 18 1
= − + + = = (1)
12 3 18 36 36 2
6
Hence
1
AG = 1 − (2)
2
too. To answer the options, we need not only the values of AG and AR ,
but also a graph of f (x). By direct calculation,
17
f (0) = , f (1) = 13/36 (3)
36
Also, f ′ (x) = x2 − 2x + 95 which has two roots, 13 and 53 , of which only
the first lies in [0, 1]. f ( 13 ) = 324
181
again by a straight calculation. So it
is the maximum for f (x) on [0, 1]. f (x) increases from 17 36
to 181
324
and
13
thereafter decreases to 36 , which is the minimum of f (x) on [0, 1].
Coming to the options, all of them are about h ∈ [ 14 , 32 ]. We have
shown by dotted lines the extreme possibilities: h = 14 and h = 23 and
181
also Lh for h = 324 .
We are now ready to tackle the options one-by-one. Note that most
of the green region lies above the line L181/324 . The entire region above
L3/4 is green and its area is 1 × (1 − 34 ) = 41 = 12 AG . So for any h < 34
the area of the green region below Lh is smaller than that above. Hence
(A) is false. However, for (B), the bisection is of the red region and the
line L1/4 does it (just as the line L3/4 bisects the green region). Hence
(B) is true.
In (B) we explicitly gave a value of h ∈ [ 14 , 32 ] which worked. This
may not always be easy or even possible. So we resort to what are called
existence theorems. These are theorems which assert the existence of
something without giving any explicit formula for it. Rolle’s theorem,
or more generally, Lagrange MVT, or the Intermediate Value Property
of continuous functions (used in the solution to Q.1 above) are well-
known examples of existence theorems. In (C) we have to compare
the green area above Lh (abbreviation for the area of the portion of
the green region above Lh ) with the red area below Lh . Denote these,
respectively, by G+ −
h and Rh . (We can also define their complements
Gh and Rh . They are related by G+
− + − 1 +
h + Gh = 2 and Rh + Rh = 2 .)
− 1

Clearly G+ −
h decreases as h increases and Rh increases as h does.

The crucial point to observe is that G+


h is a continuous function of
h. This follows because for any h and ∆h, G+ +
h − Gh+∆h is the green

7
area between the strip bounded by Lh and Lh+∆h and hence

|G+ +
h+∆h − Gh | ≤ |∆h | (4)

since the strip is a rectangle with sides 1 and |∆h|. Similarly, Rh− is a
continuous function of h on [0, 1]. (Actually, these are understatements.
For those who know it, we have shown that G+ −
h and Rh satisfy what is
called a Lipschitz condition which trivially implies continuity.)
Coming to (C), it asks whether there exists some h ∈ [ 14 , 23 ] for which

G+ −
h − Rh = 0 (5)

Since 13
36
is the minimum of f (x) on [0, 1], the entire green region lies
above L13/36 . Hence G+ 1 13
h = 2 for h = 36 . But the red area below L13/36
13
is only 36 . Hence

G+ −
h − Rh > 0 (6)
13 181
for h = 36 . By a similar reasoning and using that 324
is the maximum
of f (x) over [0, 1] we get that

G+ −
h − Rh < 0 (7)

for h = 181
324
. From (6), (7) and the IVP, there exists some h ∈ [ 36 13 181
, 324 ]
+ − 13 181 1 2
for which Gh − Rh = 0. Since [ 36 , 324 ] ⊂ [ 4 , 3 ], we get that (C) is true.
Finally, with the notations introduced above, (D) amounts to
showing that

Rh+ = G−
h (8)

for some h ∈ [ 41 , 32 ]. This can be proved by a similar argument as for


(C). But that is not necessary. Since Rh+ = 21 − Rh− and G− 1 +
h = 2 − Gh ,
truth of (D) follows from that of (C).

Basically, a good problem on theoretical calculus based on IVP. But


the computational part is also considerable and a mistake in it could
affect the answer since in (C) and (D), it was crucially needed that the
minimum and the maximum of f (x) on [0, 1] both lie in [ 14 , 32 ].

8
SECTION - 2
This section contains FOUR questions. Each question has FOUR options
of which ONLY ONE is correct. There are 4 marks for choosing only the
correct option, no marks if no option is chosen and -1 mark in all other cases.

Q.4 Let
h f : (0, 1) −→ IR be the function defined as f (x) = n if x ∈
1
, 1 where n ∈ IN. Let g : (0, 1) −→ IR be a function such that
n+1 ns
Z x 1−t √
dt < g(x) < 2 x for all x ∈ (0, 1). Then lim f (x)g(x)
x2 t x→0

(A) does NOT exist


(B) is equal to 1
(C) is equal to 2
(D) is equal to 3

Answer and Comments: (C). Nothing is given about the function


g(x) other that two bounds for it. The only way we can find from this
the given limit is to use these bounds to get ‘suitable’ upper and lower
bounds for f (x)g(x) so that the Squeeze Theorem (more popularly
called the Sandwich Theorem) can be applied.
In notations, we seek two non-negative functions U(x) and L(x)
such that

L(x) ≤ f (x) ≤ U(x) (1)

for all x ∈ (0, 1). Since g(x) is also non-negative everywhere this would
imply that

L(x)g(x) ≤ f (x)g(x) ≤ U(x)g(x) (2)

and further, using the given bounds for g(x),


s
Z x 1−t √
dtL(x) ≤ f (x)g(x) ≤ 2 xU(x) (3)
x2 t
Suitability of the q
bounds U(x) and L(x) means that as x → 0, both
√ Rx
2 xU(x) and x2 1−t t
dtL(x) should tend to the same limit. If such

9
bounds do not exist then lim f (x)g(x) may still exist but we will not
x→0
be able to find it from the data.
1
Suitability of U(x) simply means that it should be of the order of √ .
x
1
As a simple choice we may take U(x) to be √ itself. Of course, we
x
must, in the first place, ensure that it is an upper bound for f (x) for
all x ∈ [0, 1]. But that is easy. From the very definition,

f (x) = n (4)
1
where n is such that n+1 ≤ x < n1 . This implies that nx < 1 and hence

nx < 1. Therefore,
√ 1
f (x) = n≤ √ (5)
x

Finding a lower bound L(x) for f (x) is easy. Proving its suitability is
not so easy. We have (n + 1)x ≥ 1. Hence

nx ≥ 1 − x (6)

which gives
s
√ 1−x
f (x) = n≥ (7)
x
So we may take
s
1−x
L(x) = (8)
x
To sprove that L(x) is suitable we must now prove not only that
Z x 1−t
dtL(x) tends to some limit as x → 0, we must prove that it
x2 t
√ 1
tends to 2, as that is the limit of 2 xU(x), if we take U(x) as √ . For
x
this, we first evaluate the integral. Using the substitution t = sin2 u,

10
the indefinite integral comes as
s
Z Z
1−t
dt = 2 cot u sin u cos udu
t Z
= 2 cos2 udu
= (1 + cos 2u)du
= u + sin u cos u (9)

Hence
s
Z x 1−t √ √ √ √
dt = sin−1 x + x 1 − x − sin−1 x − x 1 − x2 (10)
x2 t
s
Z
1−t x
From (8) and (10), lim L(x) dt is the same as
√ √
x→0
√ √ x 2 t √
1 − x(sin−1 x + x 1 − x − sin−1 x − x 1 − x2 )
lim √
x→0 x

√ sin−1 x sin−1 x
As x → 0+ , 1 − x and √ both tend to 1. But √ tends to
x√ x
0 since x tends to 0 faster than x. Therefore the first two terms inside
the parentheses contribute 1 each to the limit
s while the last two terms
Z x
1−t
contribute nothing. All put together, L(x) → 2 as x → 0+ .
x2 t
Therefore by the Sandwich Theorem, lim f (x)g(x) = 2.
x→0

s √
Z x 1−t
1−x √
In finding lim √ dt ×
, we may drop 1 − x (which
x→0 x2 tx
0
tends to 1). Thereafter, the expression reduces to an indeterminate
0
form and its limit may be computed using L’ôpital’s rule, using the
fundamental
s theorem
s of calculus
s to get the derivative of the numerator
1−x 1−x 2 1−x √
as − 2x = − 2 1 − x2 . The derivative of the
x x2 x
√ 1
denominator x is √ . Therefore the limit of the L.H.S. of (3) equals
√ 2 x
√ 1−x √
lim 2 x( √ − 2 1 − x2 ) which comes out as 2. This obviates the
x→0 x

11
s
Z
1−t
x
need to evaluate the integral dt.
x2 t
A delicate problem where squeeze is the only way and to apply it
one needs an upper bound U(x) as well as a lower bound L(x) on f (x)
valid for all points of (0, 1) (or at least for all points in some deleted
neighbourhood of 0) and both of which are of the order of √1x . Although
both are easy to obtain, taking U(x) as √1x is more obvious. Thereafter,
1
a hasty argument, using constancy of f on intervals of the form [ n+1 , n1 )
may give the wrong impression that √1x is also a lower bound on f (x).
And the answer will come even with this mistake. Were it not for this
possibility of rewarding the unscrupulous candidates, it is a very good
problem.

Q.5 Let Q be the cube with the set of vertices (x1 , x2 , x3 ) : x1 , x2 , x3 ∈


{0, 1}. Let F be the set of all twelve lines containing the diagonals of
the six faces of the cube Q . Let S be the set of all four lines containing
the main diagonals of the cube Q ; for instance, the line passing through
the vertices (0, 0, 0) and (1, 1, 1) is in S. For lines ℓ1 and ℓ2 , let d(ℓ1 , ℓ2 )
denote the shortest distance between them. Then the maximum value
of d(ℓ1 , ℓ2 ) as ℓ1 varies over F and ℓ2 varies over S , is
1 1 1 1
(A) √ (B) √ (C) √ (D) √
6 8 3 12

Answer and Comments: (A). A cube has 8 vertices. Call (0, 0, 0)


as O, (1, 0, 0) as A, (0, 1, 0) as B and (0, 0, 1) as C. Through each one
of these there is a diagonal of the cube and it is customary to denote
their other extremities by putting dashes over the labels of the starting
points. Thus OO ′ is a main diagonal where O ′ = √ (1, 1, 1). So is BB


where B = (1, 0, 1) and so on. Each has length 3. These four main
diagonals form the set S in the statement of the problem.
The cube also has 6 square faces, all congruent to each other other.
Each has two diagonals. For example, the top face (with vertices
O ′, A′ , C, B ′ has A′ B ′ and CO ′ as its diagonals.
√ Let us call these 12
diagonals as facial diagonals. Each has length 2. They constitute the
set F in the statement of the problem.

12
C (0,0,1)
A’(0,1,1)

B’
(1,0,1) O’ (1,1,1)
O B (0,1,0)
(0,0,0)

C’
A (1,0,0)

The problem asks the maximum distance between pairs of diagonals, of


which one is a main diagonal and the other a facial one. By symmetry,
we may assume that the main diagonal is OO ′. Now the 12 facial diag-
onals can be classified into two classes depending on their relationship
with OO ′. There are six facial diagonals of which three pass through
O and the other three pass through O ′ . Clearly the shortest distance
of any of these from OO ′ is 0.
The other six facial diagonals are skew w.r.t. OO ′, that is, the lines
containing them do not meet the line containing OO ′. For example,
AB, A′ B ′ and AC are such facial diagonals. By symmetry the shortest
distance between OO ′ and any of these six diagonals is the same. So,
without loss of generality, we calculate the shortest distance between
the line OO ′ and the line AB.
There is a well-known formula for the shortest distance between any
two skew lines. It is the projection onto their common normal of any
segment having one point on each line. A vector parallel to OO ′ is
î + ĵ + k̂ while a vector along AB is −î + ĵ. Hence a common normal
~n to both is

~n = (î + ĵ + k̂) × (−î + ĵ)

13
î ĵ k̂
= 1 1 1
−1 1 0
= −î − ĵ + 2k̂ (1)

Taking O as a point on the main diagonal and A as a point on the


−→
short diagonal, the vector OA equals î. Hence the shortest distance,
say d, between OO ′ and AB equals

|î · (−î − ĵ + 2k̂)|


d =
| − î − ĵ + 2k̂|
1
= √ (2)
6

So the maximum shortest distance between any main diagonal and any
1
facial diagonal is √ .
6

A simple problem once it is realised that by symmetry, the shortest


distance between a skew pair of a main diagonal and a facial diagonal is
the same for any such pair. Hence it is far from a typical maximisation
problem in calculus!
( )
x2 y 2
Q.6 Let X = (x, y) ∈ ZZ × ZZ : + < 1 and y 2 < 5x . Three dis-
8 20
tinct points P, Q, R are randomly chosen from X. The the probability
that P, Q and R form a triangle whose area is a positive integer, is
71 73 79 83
(A) (B) (C) (D)
220 220 220 220

Answer and Comments: (B). Verbally, the set X consists of all


points with integer coordinates which lie inside the ellipse, say E and
between the upper and the lower arcs of the parabola. The second
restriction rules out points for which√the x-coordinate is negative. The
first restriction implies that |x| ≤ 8. Both the restrictions, put to-
gether imply that in order for a point (x, y) to lie in X, x can take only
two values, 1 and 2.

14
For x = 1, (1, y) will lie inside the ellipse if and only if
1 35
y 2 < 20(1 − ) = (1)
8 2
while it will lie between the upper and the lower branch of the parabola
if and only if

y2 < 5 (2)

As the second restriction is stronger than the first, we enumerate only


those points (1, y) which satisfy (2). They constitute a subset, say X1 ,
of X given by

X1 = {(1, 0), (1, 1), (1, −1), (1, 2), (1, −2)} (3)

When x = 2, the restriction y 2 < 10 is the same as the restriction


y 2 < 20(1 − 12 ) = 10. So we get another subset X2 of X given by

X2 = {(2, 0), (2, ±1), (2, ±2), (2, ±3)} (4)

Since X = X1 ∪ X2 , it has 5 + 7 = 12 elements. Three (distinct) points


P, Q, R can be chosen from it randomly in
!
12
= 220 (5)
3
ways. For these three points to form a triangle, two must come from
X1 and 1 from X2 or vice versa. The resulting triangle will always have
width 1. For the area to be an integer, the height, i.e. the vertical
distance between the two points lying in the same Xi must be even.
Put differently, the y-coordinates of these two points must have the
same parity, i.e. either both are even or both odd.
   
When two points are from X1 , this happens in 32 + 22 = 4 ways. The
third vertex is from X2 and can be chosen in 7 ways. So there are 28
triangles of this type.
Similarly, when two vertices
   of the triangle are from X2 (and the third
3 4
from X1 ), there are ( 2 + 2 ) × 5 = (3 + 6) × 5 = 45 triangles.
Together there are 28 + 45 = 73 favourable cases. So the desired
73
probability is .
220
15
A simple but interesting combination of coordinate geometry and
counting problems.

Q.7 Let P be a point on the parabola y 2 = 4ax, where a > 0 . The normal
to the parabola at P meets the x-axis at a point Q . The area of the
triangle P F Q , where F is the focus of the parabola, is 120. If the
slope m of the normal and a are both positive integers, then the pair
(a, m) is

(A) (2,3) (B) (1,3) (C) (2, 4) (D) (3, 4)

Answer and Comments: (A). It is convenient to take the parametric


representation of the parabola as

x = at2 , y = 2at (1)

One of the advantages of this representation is that the slope of the


normal to the parabola at the point P = (at2 , 2at) is simply −t.
The equation of the normal at P is

y − 2at = −t(x − at2 ) (2)

It meets the x-axis when y = 0 and hence x = 2a + at2 . So the point Q


is (at2 + 2a, 0). The area of △P F Q is easy to find since both F = (a, 0)
and Q lie on the x-axis. So the base is at2 + a and the height is the
y-coordinate of P , i.e. 2at.
Therefore

area of △P F Q = |a2 t(t2 + 1)| (3)

This is given to equal 120. Also as the answer wanted is in terms of m


rather than t, we replace t by −m to get

|a2 m(m2 + 1)| = 120 (4)

It would have been an interesting problem to find all integral solutions


(a, m) of this equation. But we are only asked which of the four options
is a solution. This can be done by direct substitution. Luckily, the very
first option where a = 2 and m = 3 works. As only one option is correct,
we need not try the others. So (A) is the correct option.

16
An absolutely simple problem. The geometry part is simple, once
you realise how the slope of the normal is related to the parameter.
The number theory part consists merely of a verification rather than
solving an integral equation.

SECTION - 3
This section contains SIX questions. The answer to each question is a
non-negative integer. There are 4 points ONLY for a correct answer and 0
in all other cases.
 
π π
Q.8 Let tan−1 (x) ∈ − , , for x ∈ IR. Then the number of real so-
2 2 q √
lutions of the equation 1 + cos(2x) = 2 tan−1 (tan x) in the set
     
3π π π π π 3π
− ,− ∪ − , ∪ , is equal to
2 2 2 2 2 2

Answer and Comments: 3. As the set consists of three disjoint


intervals, it suffices to find the number of solutions in each and add.
Because 1 + cos(2x) = 2 cos2 x, the given equation simplifies to

| cos x| = tan−1 (tan x) (1)

For x ∈ (0, π2 ), tan−1 (tan x) = x. Also cos x ≥ 0. So the equation


reduces to cos x = x, or equivalently, to cos x − x = 0. The L.H.S. is
positive at 0 and negative at points near π2 . So by the Intermediate
Value Property, it has at least one solution in [0, π2 ). At the same
time, the L.H.S. is a strictly decreasing function of x as we see from its
derivative. Hence it cannot have more than one solution. Put together,
there is exactly one solution in [0, π2 ). However, in (− π2 , 0] there is no
solution because | cos x| is positive while tan−1 (tan x) = x is negative.
Summing up, the given equation has exactly one solution in (− π2 , π2 ).
The | cos x| function is periodic with period π. So its values over the
middle interval repeat over the other two intervals. tan x is also periodic
with period π. It is, therefore, tempting to think that the other two
intervals also have exactly one solution. But there is a catch. Even
though tan x is periodic, because of the stipulation that tan−1 always

17
takes values in (− π2 , π2 ), we no longer have tan−1 (tan x) = x. But again,
because of periodicity, for x ∈ (− π2 , 3π 2
) we have

tan−1 (tan x) = x + π (2)

Hence if x0 is a solution of (1) in (− π2 , π2 ), then its translate x0 + π will


be a solution of (1) in ( π2 , 3π
2
). Conversely, if x1 is a solution of (1) in
the latter, then its translate x1 − π will be a solution in the former. So
the number of solutions of (1) is the same in both, i.e. 1 By a similar
reasoning, there is one solution in the third interval (− 3π 2
, − π2 ).
So, in all the given equation has three solutions.

A simple but good question on inverse trigonometric functions. It


would have been better as a question where reasoning has to be given
because some candidates may conclude the correct answer hastily from
the mere periodicity of the | cos x| function and of the tan x function.
The problem can also be tackled graphically. The graph of y = | cos x|
consists of recurring arcs. Once the idea that the graph of y = tan−1 (tan x)
consists of translates of line segments with slope 1, the solution is in-
stantaneous. We have preferred to give an analytic solution.

Q.9 Let n be a natural number and f : [0, 1] −→ IR be the function defined


by 
1
 n(1 − 2nx) if 0 ≤ x ≤ 2n







 1 3


 2n(2nx − 1) if 2n
≤x≤ 4n
f (x) =
 3 1



 4n(1 − nx) if 4n
≤x≤ n





 n 1

n−1
(nx − 1) if n
≤x≤1

If n is such that the area of the region bounded by the curves x =


0, x = 1, y = 0 and y = f (x) is 4 , then the maximum value of the
function f is

Answer and Comments: 8. The value of n is not given. But no


matter what it is, the function is piecewise linear and its graph consists
of four straight line segments as shown below.

18
y

n (1,n)

3/4n
.
O 1/2n 1/n (1,0) x

We are given that the area below the graph of f is 4. As the graph
consists of straight line segments, this area is the sum of the areas of
the three triangles. Adding, we get
 
n 1 1 1
+ + (1 − ) = 4 (1)
2 2n 2n n
On simplifying, this leads to n = 8. This is also the maximum value of
f (x).

A simple problem once it is understood what it is asking.


r
z }| {
Q.10 Let 7 5 . . . 5 7 denote the (r + 2) digit number where the first and the
last digits are 7 and the remaining r digits are 5. Consider the sum
99
98 z }| {
z }| { 75...57 + m
S = 77 + 757 + 7557 + 7 5 . . . 5 7. If S = , where m and
n
n are natural numbers less than 3000, then the value of m + n is

Answer and Comments: 1219. Call the given number as ar . Then


a0 = 77 and the sum S equals
98
X
S= ar (1)
r=0

19
We are then asked to determine the integers m and n for which
a99 + m
S= (2)
n
The problem may not have a unique solution because (2) defines a single
linear equation nS = a99 + m in the two unknowns m and n. If we
allow m, n to take real values, there will be infinitely many solutions.
Because of the stipulation on m and n (viz., that they are positive
integers less than 3000) the solution will be unique and we have to find
it.
Clearly S is a huge sum. The stipulation that m, n are less than 3000
can also be looked at as a way of telling us that the dominating part
a99
of the second expression for S is . There is, of course, nothing
n
special about 98 and 99. They are consecutive positive integers and
large enough to preclude the possibility of solving the problem by brute
force summation.
So the crucial question is for which value of the positive integer n,
a99
a0 + a1 + . . . + a98 nearly equals . The most familiar example where
n
the sum of the first 99 terms of some series nearly equals some constant
multiple of its 100-th term is when the series is a geometric series.
Specifically, if k 6= 1 is the common ratio of a G.P. {cn }n≥0 , then
c0
c0 + c1 + . . . + c98 = (k 99 − 1) (3)
k−1
and, more generally, for any positive integer p,
c0
c0 + c1 + c2 + . . . + cp−1 = (k p − 1) (4)
k−1

A geometric progression with common ratio k is nothing but a solution


of a simple recurrence relation viz.

cr = kcr−1 (5)

for r ≥ 1.
The proof of (4) using (5) is often given as a standard example of a
proof by mathematical induction. But like all proofs by induction it

20
has the drawback that you have to know or at least guess the answer
beforehand. You can only verify it.
What if we have to arrive at (4) from (5)? Then there is an ingenious
way by manipulating the sum, say Sp on the left of (4). We rewrite (4)
except that we now replace the R.H.S. by Sp which is to be found.

c0 + c1 + c2 + . . . + cp−2 + cp−1 = Sp (6)

If we multiply this by k we get

kc0 + kc1 + . . . + kcp−2 + kcp−1 = kSp (7)

If we subtract (6) from (7), we see that because of (5), most of the
terms get canceled and we are left with

(k − 1)Sp = kcp−1 − c0 = cp − c0 = c0 k p − c0 = c0 (k p − 1) (8)

which proves (4).


(5) is an example of what is called a linear recurrence relation of
order 1. (Order indicates how many preceding terms determine a given
term. Thus the well-known Fibonacci relation Fr = Fr−1 + Fr−2 is of
order 2.) The alternate derivation of (4) amounts to saying that we
have solved the recurrence relation (5), as opposed to merely verifying
a solution of it. The theory of such relations runs parallel to the theory
of linear differential equations. This is not a coincidence because it is
indeed possible to convert linear recurrence relations for sequences to
linear differential equations using what are called exponential gen-
erating functions of sequences. But discussing it will take us too
far afield. (Those interested may see Exercise (3.32) of Chapter 6 of
Foundations of Discrete Mathematics by the author.)
Returning to our problem, it is tempting to guess the integers m and n
a2 + m
by trying some initial special cases such as a0 + a1 = and then
n
prove by induction that for every positive integer p
ap + m
Sp = (9)
n
where Sp = a0 + a1 + . . . + ap−1 . But the trouble is that just as m, n
are not unique as observed above, the values we get may vary with p.

21
So, we must go through the alternate approach, viz. to first write a
suitable recurrence relation and then to solve it.
It is a little awkward to write a linear recurrence relation for the given
sequence {ar }r≥0 . But if we remove the 7 in the unit place and define
r
z }| {
br = 7 5 . . . 5 (10)

then ar can be easily recovered from br as

ar = 10br + 7 (11)

The advantage {br }r≥0 has over {ar }r≥0 is that it satisfies a simple
recurrence relation, viz.

br = 10br−1 + 5 (12)

with b0 = 7. Had the last term 5 on the R.H.S. been missing, this
would be a recurrence relation like (5) and will make {br }r≥0 into a
geometric progression with common ratio 10. As it now stands, it is a
non-homogeneous linear recurrence relation of order 1, while (5) is
a homogeneous linear recurrence relation of order 1.
Despite this difference, let us see if the method used in proving (4) can
be pushed through with suitable modifications. We let
98
X
T = br = b0 + b1 + b2 + . . . + b98 (13)
r=0

The role of the coefficient k is played by 10 here. So, if we multiply T


by 10 we get

10T = 10b0 + 10b1 + 10b2 + . . . + 10b98 (14)

Now if we subtract (13) from (14) (analogously to subtracting (6) from


(7)), the terms will not cancel, because, for example, b1 does not equal
10b0 but rather 10b0 + 5. So to make such cancellations possible, let us
add 5 to each of the 99 terms on the right of (14) and then subtract
(13). Since 99 × 5 = 495, we get

9T + 495 = 10b98 + 5 − b0 = b99 − 7 (15)

22
and hence
b99 − 502
T = (16)
9
We are almost done. Now it is merely a matter of relating T with
S and b99 with a99 . This is done through (11). Applying it for r =
0, 1, 2, . . . , 98 and adding, we get

S = a0 + a1 + . . . + a98
= 10T + 99 × 7
= 10T + 693
b99 − 502
= 10 + 693 (17)
9
10b99 − 5020 + 6237
=
9
10b99 + 1217
=
9
a99 + 1210
= (18)
9
where, in (17) we have used (16) and in (18) we have used (11).
Thus finally we get that m = 1210 and n = 9 and, therefore, m + n =
1219.

A good problem. Those who have solved linear recurrence relations


earlier will have an easier time in recognising that taking 10S and then
subtracting S from it is the key step. In fact, once this idea strikes,
the solution can be completed more efficiently bypassing the auxiliary
sequence {br }r≥0 completely. We note that for every r ≥ 1, the numbers
10ar−1 and ar differ only in their last two digits. 10ar−1 ends with 70
as the last two digits while ar ends with 57. Since 70 − 57 = 13, we
have

10ar−1 = ar + 13 (19)

for all r ≥ 1. Summing over from r = 1 to r = 99 gives


99
X 99
X
10 ar−1 = ar + (13 × 99) (20)
r=1 r=1

23
The L.H.S. is precisely 10S. The sum on the R.H.S. is S except that
a0 is missing and a99 is added. Therefore,

10S = S + a99 − a0 + 1287 (21)

Since a0 = 77, we get

9S = a99 + 1210 (22)

or,
a99 + 1210
S= (23)
9
exactly as before.
Strictly speaking, neither method gives a complete answer. m =
1210 and n = 9 is one possible solution of (2). But as observed earlier,
(2) does not determine m, n uniquely. What if there are other integral
solutions? To investigate this, we analyse
a99 + m a99 + 1210
= (24)
n 9
which simplifies to

9m = (n − 9)a99 + 1210n (25)

This means that the R.H.S. is divisible by 9. We reduce it modulo 9.


1210 is congruent to 4 modulo 9. In a99 two digits are 7 each while
the remaining 99 are 5 each. As 99 is a multiple of 9, the sum of the
digits of a99 is congruent to 14 and hence to 5 modulo 9. Hence a99
is itself congruent to 5 modulo 9. All put together, the R.H.S. of (25)
is congruent to 5n + 4n = 9n modulo 9. Hence it is divisible by 9, no
matter what n is and will give an integral value of m. To eliminate
the unwanted possibilities, not that a99 is a huge number. So, unless
n = 9, the value of m we get will lie well outside the range 1 to 3000.
Therefore the only permissible solution of (2) is m = 1210 and n = 9.
It is unlikely that any candidate will be scrupulous enough (read
stupid enough) to weed out other solutions. But the papersetters are to
be commended for giving the stipulation that m ≤ 3000 which ensures
uniqueness of the solution.

24
1967 + 1686i sin θ
Q.11 Let A = { : θ ∈ IR}. If A contains exactly one
7 − 3i cos θ
positive integer n, then the value of n is

Answer and Comments: 281. The set A is the range of a function


f from IR to C
| defined by

1967 + 1686i sin θ


f (θ) = (1)
7 − 3i cos θ
We are given that this function assumes only one positive integer as a
value and we are asked to identify it.
For a methodical solution, we first identify those values of θ for which
f (θ) is real and then further those among them when f (θ) is an integer,
and finally, when it is a positive integer.
For the first task, we need to set the imaginary part of f (θ) to 0. First
we simplify f (θ). As 281 is a common factor of 1967 and 1686. So,
281(7 + 6i sin θ)
f (θ) =
7 − 3i cos θ
(7 + 6i cos θ)(7 + 3i cos θ)
= 281 ×
49 + 9 cos2 θ
49 − 18 sin θ cos θ + i(21 cos θ + 42 sin θ)
= 281 × (2)
49 + 9 cos2 θ
Hence f (θ) will be real if and only if 21 cos θ + 42 sin θ = 0. Again we
are lucky that 21 is a common factor. So, f (θ) is real if and only if

cos θ + 2 sin θ = 0 (3)

This is an easy trigonometric equation whose solutions are those θ for


which tan θ = − 12 . Our interest, however, is not so much to find all
solutions of this as to find the values of sin θ cos θ and cos2 θ from it,
which are needed to find the real part of f (θ).
tan θ = − 12 gives two possibilities, either sin θ = √
−1
5
and cos θ = √25 , or,
sin θ = √15 and cos θ = √−2
5
. But no matter which possibility holds, we
get
−2 4
sin θ cos θ = and cos2 θ = (4)
5 5
25
Substituting these along with (3) into (2) we get that when f (θ) is real,
49 + 36
5
its value is 281 × 36 . We are again lucky that this is simply 281.
49 + 5
Hence when f (θ) is real, it is already a positive integer 281. So, this is
the only positive integer in the set A.

The particular formatting of the question makes a sneaky solution


(given by Vansh Gaur) possible. Simply set f (θ) = n and rewrite it as
1967 + 1686i sin θ = 7n − 3ni cos θ (5)
Equating the real parts of both the sides gives 1967 = 7n and hence
n = 281. Strictly speaking, for n to qualify to belong to A, one must
show that there is some θ for which the imaginary parts are equal too
for some θ. But when we are already given that there is exactly one
value of n in A, this is unnecessary. A better question to preclude such
a sneaky answer would have been to ask how many integers are in the
set A, the alternatives being given as, say, (A) none, (B) one, (C) more
than one but finitely many and (D) infinitely many.

Q.12 Let P be the plane 3x + 2y + 3z = 16 and let S = {αî + β ĵ + γ k̂ :
7
α2 + β 2 + γ 2 = 1 and the distance of (α, β, γ) from the plane P is }.
2
Let ~u, ~v and w
~ be three distinct vectors in S such that |~u −~v | = |~v −w|
~ =
|w−~
~ u|. Let V be the volume of the parallelepiped determined by vectors
80
~u, ~v and w.~ Then the values of √ V is
3

Answer and Comments: 45. Let O be the origin and T be the unit
sphere centred at O. The perpendicular distance of O from the plane
P is √ √ 16
2 2 2
= √1616 = 4. As 4 exceeds the radius of the sphere,
( 3) +2 +3
the entire sphere T lies on one side of P . There are two planes parallel
to P at a distance 72 from it. They lie on opposite sides of P . Let Q be
the one which lies on the same side of P as the sphere T does. Then
the perpendicular distance of O from Q is 4 − 27 = 12 and the set S is
T ∩ Q, i.e. the intersection of T and Q.
With this preamble, we are ready for the problem. We are given three
vectors ~u, ~v , w
~ in T . Then
|OU| = |OV | = |OW | = 1 (1)

26
where U, V, W are the terminal points of ~u, ~v , w.
~ The problem asks
for the volume of the parallelepiped with OU, OV and OW as three
coterminus edges. It is well-known that this volume is 6 times the
volume of the tetrahedron OUV W . So the real problem is to find the
latter.
There are many formulas for the volume of a tetrahedron. The most
basic is
h
volume = ∆ (2)
3
where ∆ is the area of its base and h is its perpendicular height.
From the preamble above, h is simply the perpendicular distance of O
from the plane Q. It is obtained by subtracting the distance between
the planes from the distance of O of from P . Thus
7 1
h=4− = (3)
2 2
Finding ∆, the area of the base triangle UV W is more involved. But
the triangle UV W is given to be equilateral. Hence it suffices to know
its circumradius, say R to find ∆.
By symmetry, the foot, say M, of the perpendicular from O to the
plane Q is also the circumcentre of the triangle UV W as shown in (a)
of the figure below. (The base triangle is shown separately in (b).)

O W
1
W
1
1/2
M
1 60
30
U 30 M U V

(a) (b)

As △OMU is right-angled at M and OU = 1 and OM = 12 , we get


s √
1 3
UM = 12 − = (4)
4 2

27
As △UV W is equilateral, UM is also the circumradius R of it. There-
fore, its area ∆ is

∆ = 2R2 sin U sin V sin W


3 3 ◦
= sin 60
2 √
33 3
=
2√ 8
9 3
= (5)
16

Hence volume of the tetrahedron OUV W√ is 13 h∆ = 61 × 9163 . Therefore,
V , the volume of the parallelepiped is 9163 . So √803 V = 45.

Finding the volume of the tetrahedron is the main theme of the


problem. Finding the volume of the parallelepiped is more of an ap-
pendage. The height of the tetrahedron is not specified directly but
has to be found from the information regarding the distances from the
plane P . The problem is more a test of the ability to weed through
some long winding data to pick up the essence.

Q.13 Let a and b be two nonzero real numbers. If the coefficient of x5 in the
70 4
 
2
expansion of ax + is equal to the
27bx
1 7
 
−5
coefficient of x in the expansion of ax − 2 , then the value of 2b
bx
is

Answer and Comments: 3. The most straightforward approach


is to calculate the coefficients of x5 and x−5 respectively, in the two
expansions.
The general term in the first expansion is
! k !
4 4−k 70k 8−3k

4 70
(ax2 )4−k = a x (1)
k 27bx k (27b)k

To get the coefficient of x5 , we set 8 − 3k = 5 which gives k = 1. Hence


70 280a3
the coefficient of x5 in the first expansion is 4a3 = .
27b 27b
28
By a similar reasoning, the coefficient of x−5 in the second expansion
35a3
comes as 4 .
b
Equating the two coefficients gives an equation in the two unknowns a
and b, viz.
280a3 35a3
= 4 (2)
27b b
This single equation cannot determine a and b in general. But in the
present case, we are given that a and b are non-zero. This reduces (2)
to
27
b3 = (3)
8
Further, b is real. The only real root of (3) is 23 . Hence 2b = 3.

An absolutely straightforward problem. But the calculations are


a little messy. A far more challenging problem would be to give two
equalities of coefficients of powers. That would give two equations in a
and b. But solving them would demand more work than can be justified
for 4 marks.

SECTION - 4
This section contains FOUR Matching List Sets. Each set has ONE Mul-
tiple Choice Question. Each set has two lists : List-I and List-II. FOUR
options are given in each Multiple Choice Question based on List-I and
List-II and ONLY ONE of these four options satisfies the condition asked
in the Multiple Choice Question. There are 3 marks if ONLY the option
corresponding to the correct combination is chosen, 0 marks if none of the
options is chosen and −1 mark in all other cases.

Q.14 Let α, β and γ be real numbers. Consider the following system of linear
equations
x + 2y + z = 7
x + αz = 11
2x − 3y + βz = γ

29
Match each entry in List-I to the correct entries in List-II.

List-I List-II
1
(P) If β = (7α − 3) and γ = 28, (1) a unique solution
2
then the system has
1
(Q) If β = (7α − 3) and γ 6= 28, (2) no solution
2
then the system has
1
(R) If β 6= (7α − 3) where α = 1 (3) infinitely many solutions
2
and γ 6= 28, then the system has
1
(S) If β 6= (7α − 3) where α = 1 (4) x = 11, y = −2 and z = 0
2
as a solution
and γ = 28, then the system has
(5) x = 15, y = 4 and z = 0
as a solution
The correct option is:

(A) (P ) → (3) (Q) → (2) (R) → (1) (S) → (4)


(B) (P ) → (3) (Q) → (2) (R) → (5) (S) → (4)
(C) (P ) → (2) (Q) → (1) (R) → (4) (S) → (5)
(D) (P ) → (2) (Q) → (1) (R) → (1) (S) → (3)

Answer and Comments: (A). We are given a system of three linear


equations in three unknowns x, y, z. The possibilities regarding how
many solutions it has depend on the values of the determinant

1 2 1
∆= 1 0 α (1)
2 −3 β

and the three determinants, say, ∆1 , ∆2 and ∆3 obtained by replacing


the first, the second
 and the third column, respectively, of ∆ by the

7
column vector  
 11 .
γ

30
Direct calculations give

∆ = 7α − 2β − 3 (2)
∆1 = 21α − 22β + 2αγ − 33 (3)
∆2 = 4β + 14α − αγ + γ − 22 (4)
and ∆3 = 56 − 2γ (5)

In (P), all four determinants vanish because of the conditions given.


As ∆ = 0, either there are infinitely many solutions or there are none
depending upon whether the system is consistent or not. To check
which possibility holds, the first two equations, say E1 and E2 , are
always mutually consistent. (If α 6= 1, setting y = 0 gives a solution.
If α = 1, y = −2 gives a solution.) Further, the third equation, say E3 ,
becomes
7 3
2x − 3y + ( α − )z = 28 (6)
2 2
By inspection (or a minor calculation) E3 is a linear combintion of
E1 and E2 , specifically, E3 = − 32 E1 + 27 E2 . Hence whenever E1 , E2
are satisfied, so is E3 . That is, the system of the three equations is
consistent and has infinitely many solutions. Thus (P) −→ (3). In
(Q), however, ∆ vanishes but ∆3 does not. Hence there is no solution.
So (Q) −→ (2).
This rules out (C) and (D). In both (A) and (B), (S) is given to imply
(4). So we need not check its truth. (If it is false no option is correct.)
So the answer depends now on (R). In (R) ∆ 6= 0 by the first part
of the hypothesis and so regardless of the second part, the system has
a unique solution. So, even without checking for (S), (A) is the right
option.

A straightforward but a bit tedious problem on systems of linear


equations. The duplication of implications in some pairs of options
and the elimination of an option as soon as it is known to contain
at least one incorrect implication saves some work. (Q.16 in Paper
1 of JEE 2022 Advanced was an extreme example where the correct
option could be identified without doing hardly any computations. The
present problem is not so flagrantly sneaky.)

31
Q.15 Consider the given data with frequency distribution

xi 3 8 11 10 5 4

fi 5 2 3 2 4 4

Match each entry in List-I to the correct entries in List-II.

List-I List-II
(P) The mean of the above data is (1) 2.5
(Q) The median of the above data is (2) 5
(R) The mean deviation about the mean (3) 6
of the above data is
(S) The mean deviation about the median (4) 2.7
of the above data is
(5) 2.4
The correct option is:

(A) (P ) → (3) (Q) → (2) (R) → (4) (S) → (5)


(B) (P ) → (3) (Q) → (2) (R) → (1) (S) → (5)
(C) (P ) → (2) (Q) → (3) (R) → (4) (S) → (1)
(D) (P ) → (3) (Q) → (3) (R) → (5) (S) → (5)

Answer and Comments: (A). Straightforward calculation gives


X
fi = 5 + 2 + 3 + 2 + 4 + 4 = 20 (1)
X
and fi xi = = 15 + 16 + 33 + 20 + 20 + 16 = 120 (2)
120
Hence the mean is 20
= 6. So (P) −→ (3).
As the sample size 20 is even the median is the average of the 10-th and
the 11-th observation arranged in increasing order. In their ascending
order, the observations are 3, 4, 5, 8, 10 and 12 with respective fre-
quencies 5, 4, 4, 2, 2 and 3. Since 5 + 4 < 10 < 11 < 5 + 4 + 4, both
the 10-th and the 11-th observations are 5. So their mean is 5. So the
median is 5+5
2
= 5. Hence (Q) −→ (2).

32
As in the last problem, the correct option is now already narrowed
down to (A) and (B). Since both are given to match (S) with (5), there
is no point in checking it. But we check (R).
As the mean is 6, the mean deviation about the mean is
P
fi |xi − 6|
P (3)
fi
The denominator is already calculated as 20. The numerator equals
5 × 3 + 2 × 2 + 3 × 5 + 2 × 4 + 4 × 1 + 4 × 2 = 15 + 4 + 15 + 8 + 4 + 8 = 54 (4)
54
Hence the mean deviation about the mean is 20
= 2.7.
Thus (R) −→ (4). Hence (A) is the correct option.

Statistics has been brought back into the JEE after a very long
gap. The expected depth is yet to be settled. For example, it is not
clear if well-known inequalities like the Chebychev inequality are to be
included. Perhaps this is the reason the papersetters this year have
to stuck to very elementary concepts which require nothing more than
arithmetic.
Q.16 Let ℓ1 and ℓ2 be the lines r1 = λ(î + ĵ + k̂) and r2 = (ĵ − k̂) + µ(î + k̂),
respectively. Let X be the set of all the planes H that contain the
line ℓ1 . For a plane H, let d(H) denote the smallest possible distance
between the points of ℓ2 and H . Let H0 be a plane in X for which
d(H0 ) is the maximum value of d(H) as H varies over all planes in X .

Match each entry in List-I to the correct entries in List-II.


List-I List-II

(P) The value of d(H0 ) is (1) 3
1
(Q) The distance of the point (0, 1, 2) from H0 is (2) √
3
(R) The distance of origin from H0 is (3) √
0
(S) The distance of origin from the point of (4) 2
intersection of the planes y = z, x = 1 and H0 is
1
(5) √
2
The correct option is:

33
(A) (P ) → (2) (Q) → (4) (R) → (5) (S) → (1)
(B) (P ) → (5) (Q) → (4) (R) → (3) (S) → (1)
(C) (P ) → (2) (Q) → (1) (R) → (3) (S) → (2)
(D) (P ) → (5) (Q) → (1) (R) → (4) (S) → (2)

Answer and Comments: (B). All the statements in List-I involve


the plane H0 . So we must begin by identifying it. It is a member of
the one parameter family X of all planes containing the line ℓ1 . So let
us first get hold of the general equation of this family. This is more
convenient to obtain if we express ℓ1 as the intersection of two planes. It
is given parametrically as the line whose points are of the form (λ, λ, λ)
for λ ∈ IR. The easiest way to express the equations of ℓ1 is

x=y=z (1)

There are infinitely many ways to express a line in space as the inter-
section of two planes. In the present case, a simple choice is to take the
two planes as y = x and y = z, i.e. as y − x = 0 and y − z = 0. Then
the general equation of a plane, say H, containing ℓ1 can be written as

λ(y − x) + (y − z) = 0 (2)

where λ is a real parameter. (That is why X is called a one parameter


family.)
The shortest distance between ℓ2 and a plane is 0 if ℓ2 intersects H.
Otherwise it is a constant as ℓ2 is parallel to H. So H0 is the plane (2)
for which k is so chosen that ℓ2 is parallel to it. This will happen if the
vector along to ℓ2 is parallel to the normal to H.
Rewriting (2) as −λx + (λ + 1)y − z = 0, a normal vector to it is

~n = −λî + (λ + 1)ĵ − k̂ (3)

A vector along ℓ2 is

î + k̂ (4)

(We get this effortlessly because ℓ2 is given parametrically.)

34
Comparing (3) and (4), for ℓ2 to be parallel to H, λ + 1 = 0, i.e.
λ = −1. Thus we have finally identified the member H0 of X for which
the distance of a point of ℓ2 from H0 is maximum. The equation of H0
is

x−z =0 (5)

(We are not, at this stage, interested in the value of d(H0 ). All that
matters is that it is positive and d(H) = 0 for every other member H of
X. Of course, we must also ensure that ℓ2 is not completely contained
in H0 , because in that case d(H0 ) would also vanish. But that is easy.
Taking µ = 0 in the given parametrisation of ℓ2 , the point (0, 1, 1) is
on ℓ2 . But it does not satisfy (5).)
Now that we know H0 , we can answer all statements in List -I. (P) asks
for the value of d(H0 ). It is the distance from H0 of any point A on ℓ2 .
We take A = (0, 1, 1). Then

|0 + 1| 1
d(H0 ) = √ =√ (6)
1+1 2
So (P) −→ (5). Hence the correct option is either (B) or (D).
For (Q), the distance of (0, 1, 2) from H0 is

|0 − 2| 2 √
√ =√ = 2 (7)
1+1 2
Hence (Q) −→ (4). As this does not happen in (D), we already know
that (B) is the correct option. Nevertheless, for the sake of complete-
ness, the distance of (0, 0, 0) from H0 is 0 since H0 passes through the
origin. Finally, the point of intersection comes out to be (1, 1, 1) by
directly solving the three √ equations (5) √
with x = 1 and y = z. Its
distance from the origin is 1 + 1 + 1 = 3. This matches with (1) in
List -II.

The real work is to identify the general equation of the family X


and then to identify H0 as a member of X parallel to the second line
ℓ2 . The calculations involved after that are trivial, repetitious, and two
of them unnecessary to answer the question.

35
Q.17 Let z be a complex number satisfying |z|3 + 2z 2 + 4z − 8 = 0, where
z denotes the complex conjugate of z. Let the imaginary part of z be
non-zero.

Match each entry in List-I to the correct entries in List-II.


List-I List-II
(P) |z|2 is equal to (1) 12
(Q) |z − z|2 is equal to (2) 4
2 2
(R) |z| + |z + z| is equal to (3) 8
(S) |z + 1|2 is equal to (4) 10
(5) 7
The correct option is:
(A) (P ) → (1) (Q) → (3) (R) → (5) (S) → (4)
(B) (P ) → (2) (Q) → (1) (R) → (3) (S) → (5)
(C) (P ) → (2) (Q) → (4) (R) → (5) (S) → (1)
(D) (P ) → (2) (Q) → (3) (R) → (5) (S) → (4)

Answer and Comments: (B). If we can get hold of the complex


number z we can calculate all expressions involving it. A complex
number z = x + iy is uniquely determined by two real variables, x
and y. We need two equations in them to determine them. WE are
given only one equation satisfied by z. But one single equation about
complex numbers is equivalent to two real equations arising by taking
real and imaginary parts.
In the present problem, we are given
|z|3 + 2z 2 + 4z − 8 = 0 (1)
Only the second and the third term have possibly non-zero imaginary
parts. Setting their sum equal to 0 gives
4xy − 4y = 0 (2)
As we are given√y 6= 0, we get x = 1. Hence z is of the form 1 + iy.
But then |z| = 1 + y 2 and z 2 = 1 − y 2 + 2iy. If we put this into (1),
we get
(1 + y)3/2 + 2(1 − y 2 + 2iy) + 4(1 − i) − 8 = 0 (3)

36
which simplifies to (1 + y 2)3/2 = 2(1 + y 2 ). Canceling (1 + y 2), we get
q
|z| = 1 + y2 = 2 (4)

That gives y = ± 3 and hence

z =1±i 3 (5)

The ambiguity about z arising out of the two possibilities does not
affect the values of the expressions asked. We already know |z|2 = 4.
Hence (P)−→ (2). As this happens in more than one options, we need
to investigate further. Direct calculations give

|z − z|2 = | ± i2 3|2 = 12 (6)
2
√ 2
|z + 1| = |2 ± i 3| = 4 + 3 = 7 (7)
and |z|2 + |z + z|2 = 4 + 4x2 = 8 (8)

So (Q) −→ (1), (R) −→ (3) and (S) −→ (5). So (B) is correct.

A straightforward problem. It would have been more interesting if


the information given about z was such that the implications could be
tested even without obtaining z.

37
PAPER 2

Contents

Section - 1 (Only One Correct Option Type) 38

Section - 2 (One or more Correct Options Type) 44

Section - 3 (Non-negative Integer Answer Type) 51

Section - 4 (Paragraph Type) 58

SECTION - 1
This section contains FOUR questions each of which has ONLY ONE
correct option. There are 3 marks if only the correct option is chosen, 0
marks if no answer is given and -1 mark in all other cases.
1
Q.1 Let f : [1, ∞) −→ IR be a differentiable function such that f (1) =
3
x3
Z x
and 3 , x ∈ [0, ∞). Let e denote the base of the
f (t)dt = xf (x) −
1 3
natural logarithm. Then the value of f (e) is
e2 + 4 loge 4 + e 4e3 e2 − 4
(A) (B) (C) (D)
3 3 3 3

Answer and Comments: (C). The given equation, viz.

x3
Z x
3 f (t)dt = xf (x) − (1)
1 3
can be converted, using the FTC, into a differential equation for y =
f (x), viz.

3f (x) = f (x) + xf ′ (x) − x2 (2)

or, in more familiar notations,


dy 2
− ( )y = x (3)
dx x
38
with the initial condition y(1) = 13 .
1
This is a linear differential equation with integrating factor e−2 ln x = x2
.
Multiplying by it, the d.e. becomes
1 dy 2 1
2
− 3 = (4)
x dx x x
y
The L.H.S. is the derivative of 2 while the R.H.S. is the derivative of
x
ln x. So the general solution is
y
= ln x + c (5)
x2
i.e.

y = x2 ln x + cx2 (6)
1
The initial condition y(1) = 3
determines c as 13 . Hence the particular
solution is
1
y = f (x) = x2 ln x + x2 (7)
3
For x = e, the R.H.S. becomes e2 ln e + 13 e2 = 34 e2 .

A routine problem on linear differential equations. Those who


remember the readymade formula for the general solution will save
some time.

Q.2 Consider an experiment of tossing a coin repeatedly until the outcomes


of two consecutive tosses are same. If the probability of a random toss
1
resulting in head is , then the probability that the experiment stops
3
with head is
1 5 4 2
(A) (B) (C) (D)
3 21 21 7

Answer and Comments: (B). This is a problem on infinitistic prob-


ability discussed in Chapter 23 and the problem is somewhat analogous
to the Main Problem of that chapter. In theory, the experiment can go
on ad infinitum if we keep on getting alternate heads and tails. This

39
can happen only in two cases : either an infinite sequence of HT or an
infinite sequence of T H. It can be shown that the probability of each
is 0 and so the experiment is sure to succeed. (This may sound para-
doxical and it is indeed so in finitistic probability where only the empty
set has probability 0. But in infinitistic probability, some non-empty
subsets can also have probability 0.)
Coming to the problem, let us say that a winning sequence is one which
ends with HH but in which no two consecutive entries are the same
till then. Depending on whether the first term is H or T , a winning
sequence falls into one of the two categories:

(i) n occurrences of HT followed by HH where n is some non-negative


integer. Thus sequences HH, HT HH, HT HT HH, . . . fall into
this category.
(ii) n occurrences of T H followed by H, for some positive integer n.
Examples are T HH, T HT HH, . . ..

The probability of HH is 19 , while the probabilities of HT and T H are


both 13 × 23 = 29 . Therefore, the probability of success in (i), say p1 , is

2 1
X
p1 = ( )n
n=0 9 9
1 1
=
9 1 − 92
1 9 1
= × = (1)
9 7 7
Similarly, p2 , the probability of success in (ii) is

X 2 1
p2 = ( )n
n=1 9 3
1 92
=
3 1 − 92
1 2 9
= × ×
3 9 7
2
= (2)
21

40
Adding (1) and (2), the desired probability equals p1 +p2 = 17 + 21
2
= 5
21
.

Problems on infinitistic probability are rarely asked in JEE. The


present problem fills the gap.
π π
Q.3 For any y ∈ IR, let cot−1 (y) ∈ (0, π) and tan−1 (y) ∈ (− , ). Then the
2 2
2
6y −1 9 − y 2π
sum of all solutions of the equation tan−1 ( 2
) + cot ( )=
9−y 6y 3
for 0 < |y| < 3, is equal to
√ √ √ √
(A) 2 3 − 3 (B) 3 − 2 3 (C) 4 3 − 6 (D) 6 − 4 3

Answer and Comments: (C). Because of the periodicity of the tan


and cot functions both with period π, their inverse functions tan−1
and cot−1 are multivalued, any two values differing by a multiple of
π. To ensure their single-valuedness, an interval of length π has to be
specified. In the present problem, tan−1 is given to lie in (− π2 , π2 ) while
cot−1 is to lie in (0, π). This has to be kept in mind while finding the
solutions of
2
6y −1 9 − y 2π
tan−1 ( 2
) + cot ( )= (1)
9−y 6y 3

for 0 < |y| < 3.


2 2
First consider the case y ∈ (0, 3). Then 9−y
6y
is positive and so cot−1 ( 9−y
6y
)
−1 6y
is the same as tan ( 9−y2 ). That reduces (1) to

6y π
tan−1 ( 2
)= (2)
9−y 3
which gives a quadratic in y, viz.
6y π √
= tan = 3 (3)
9 − y2 3
In a simplified form this is
√ 2 √
3y + 6y − 9 3 = 0 (4)

41

−6 + 144 √ √
whose roots are √ , i.e. 3 and −3 3. The latter falls outside
2 3 √
(0, 3). So (1) has only one solution in (0, 3), viz. 3.
9−y 2
Next, we consider solutions of (1) for which y ∈ (−3, 0). Here 6y
is
2
negative and so cot−1 ( 9−y6y
) lies in ( π2 , π). It no longer equals 6y
tan−1 ( 9−y 2)
6y
(which lies in (− π2 , 0)), but rather tan−1 ( 9−y 2 ) + π. Hence (1) becomes

6y 2π π
2 tan−1 ( 2
)= −π = − (5)
9−y 3 3
So, instead of (2), we now have
6y π
tan−1 ( 2
)=− (6)
9−y 6
which gives a quadratic in y, viz.
6y π 1
2
= − tan = − √ (7)
9−y 6 3

Simplifying
√ as before, this quadratic in y has roots 3 3 ± 6, of which
only 3 3 − 6 lies in (−3, 0).
√ √
Summing
√ up, (1) has two solutions, 3 and 3 3 − 6. Their sum is
4 3 − 6.

A good problem requiring careful analysis of the multivaluedness


of the inverse trigonometric functions.

Q.4 Let the position vectors of the points P, Q, R and S be ~a = î + 2ĵ −


5k̂, ~b = 3î + 6ĵ + 3k̂, ~c = 17
5
î + 16
5
ĵ + 7k̂, and d~ = 2î + ĵ + k̂ respectively.
Then which of the following statements is true?

(A) The points P, Q, R and S are NOT coplanar


~b + 2d~
(B) is the position vector of a point which divides P R internally
3
in the ratio 5 : 4
~b + 2d~
(C) is the position vector of a point which divides P R exter-
3
nally in the ratio 5 : 4

42
(D) The square of the magnitude of the vector ~b × d~ is 95.

Answer and Comments: (B). All the options are computational.


But unlike in Q.1 above they run into different directions and require
diverse methods to tackle. Because of the section formula, the calcula-
tions needed in (B) and (C) are easier than those in (A) or (D). So we
check them first. By a direct calculation, the position vector, say ~x of
the point X which divides the segment P R internally in the ratio 5 : 4
is
5~c + 4~a
~x =
9
(17î + 16ĵ + 35k̂) + 4î + 8ĵ − 20k̂
=
9
21î + 24ĵ + 15k̂
=
9
7î + 8ĵ + 5k̂
= (1)
3

On the other hand,

b + 2d 3î + 6ĵ + 3k̂ + 4î + 2ĵ + 2k̂


=
3 3
7î + 8ĵ + 5k̂
= (2)
3
From (1) and (2) we see that (B) is correct.

As only one option is given to be correct, there is no need to


check the other options. Still for the sake of completeness, (C) is false
because the same point cannot be both the internal and the external
divisor in the same ratio. Falsity of (D) can be proved by first directly
calculating ~b × d.
~

î ĵ k̂
~b × d~ = 3 6 3
2 1 1
= 3î + 3ĵ − 9k̂ (3)

43
Hence |~b × d|
~ 2 = 9 + 9 + 81 = 95 6= 99. So (D) is false.
Finally, we must prove that the points P, Q, R and S are coplanar.
This is equivalent to showing that the vectors ~b − ~a, ~c − ~a and d~ −
~a are linearly dependent. By a direct calculation, these vectors are,
respectively, 2î + 4ĵ + 8k̂, 12
5
î + 65 ĵ + 12k̂ and î − ĵ + 6k̂. For checking
linear dependence, we find the determinant, say ∆, of the coefficients,
viz.
2 4 8
12 6
∆ = 5 5
12 (4)
1 −1 6
5
Dividing the first row by 2 and multiplying the second row by 6
gives
another determinant, viz.

1 2 4
2 1 10 (5)
1 −1 6

As the second row is the sum of the other two rows, the determinant
vanishes and so does ∆. Therefore the vectors ~b − ~a, ~c − ~a and d~ − ~a
are linearly dependent and so the points P, Q, R and S are coplanar.
When only one option is given to be correct and the work needed
for the options runs into several different directions, there is an element
of luck in the order the statements are tested. In the present problem,
a candidate who successfully proves that (A) is false has wasted a lot
of time. Some intelligent thinking may guide as to which option should
be tried first. Perhaps that is what the paper-setters had in mind while
giving such diverse options.

SECTION - 2
This section contains THREE questions each of which has ONE or MORE
correct option(s). There are 3 marks if only the correct option(s) are chosen,
There is some partial credit if only some but not all correct options and no
incorrect option is chosen. There are 0 marks if no answer is given and -1
mark in all other cases.

44
Q.5 Let M = (aij ), i, j ∈ {1, 2, 3}, be the 3 × 3 matrix such that aij = 1 if
j + 1 is divisible by i, otherwise aij = 0. Then which of the following
statements is (are) true?

(A) M is invertible
   
a1 a1
   
(B) There exists a non-zero column vector  a2  such that M  a2  =
a3 a3
 
−a1
 
 −a 2 
−a3
 
0
3  
6 {0}, where 0 =  0 
(C) The set {X ∈ IR : MX = 0} =
0
(D) The matrix M − 2I is invertible, where I is the 3 × 3 identity
matrix.

Answer and Comments: (B, C). The entries of the matrix are
specified in a twisted way. For a 3 × 3 matrix (aij ), j + 1 takes the
values 2, 3 and 4. Their divisors in the set {1, 2, 3} are, respectively,
the sets {1, 2}, {1, 3} and {1, 2}. Hence for j = 1, aij equals 1 for i = 1
or 2 and 0 for i = 3. Determining the entries in the other two columns
similarly, we get
 
1 1 1
 
M = (aij ) =  1 0 1  (1)
0 1 0

Expansion of |M| w.r.t. its last row gives |M| = 0. Hence M is


singular, rendering (A) false. Singularity also implies the existence
of some non-zero column vector X such

as MX = 0. Thus (C) is
x
 
true. (An explicit column vector  y  can be obtained by noting that
z
the system x + y + z = 0, x + z = 0, y = 0 has a non-zero solution
x = 1, y = 0, z = −1.)

45
A similar argument works for (B). The equation
   
a1 −a1
   
M  a2  =  −a2  (2)
a3 −a3

is equivalent to the system a1 + a2 + a3 = −a1 , a1 + a3 = −a2 , a2 = −a3


which has a1 = 0, a2 = 1, a3 = −1 as a non-trivial solution.
(D) can be checked by directly finding |M − 2I|.
   
1 1 1 2 0 0
   
M − 2I =  1 0 1  −  0 2 0 
0 1 0 0 0 2
 
−1 1 1
 
=  1 −2 1  (3)
0 1 −2

Hence

|M − 2I| = −3 + 2 + 1 = 0 (4)

So M − 2I is not invertible. Hence (D) is false.

Sophisticated arguments can be given for the options using the


concept of eigenvalues. These are the roots of the polynomial equation
|M − xI| = 0. In the present problem

1−x 1 1
|M − xI| = 1 −x 1
0 1 −x
= (1 − x)(x2 − 1) + x + 1
= −(x + 1)x(x − 2) (5)

Thus we see that the eigenvalues are 0, −1 and 2. If λ is an eigenvalue


of a square matrix A, then M − λI is singular, and hence the system
Ax = λx has a non-trivial solution. This takes care of all the options.
Moreover we are spared of actually exhibiting the desired column vec-
tors. As it stands, there is considerable duplication of work in proving
(B) and (C). It is also not clear what is gained by specifying M in

46
a twisted manner. It would have some value if the options could be
tested simply from the divisibility conditions, without identifying M
explicitly.
1 2
  
Q.6 Let f : (0, 1) −→ IR be the function defined as f (x) = [4x] x − x − 12
4
where [x] denotes the greatest integer less than or equal to x. Then
which of the following statements is(are) true?

(A) The function f is discontinuous exactly at one point in (0, 1)


(B) There is exactly one point in (0, 1) at which the function f is
continuous but NOT differentiable
(C) The function f is NOT differentiable at more than three points in
(0, 1)
1
(D) The minimum value of the function f is − .
512

Answer and Comments: (A, B). The function f (x) is the product
of three functions. The second and the third factors are differentiable
everywhere. So, any irregularity f (x) may have has to come from the
first factor [4x]. We first identify these. They may not all be irreg-
ularities of f (x) because sometimes an irregularity of one factor gets
‘cured’ or masked if the other factor is powerful. (For example, sin x1 is
discontinuous at 0. But the product x sin x1 is continuous. When mul-
tiplied by a more powerful factor x2 , the product x2 sin x1 becomes not
just continuous but differentiable at 0. In general, if w(x) = u(x)v(x)
and u(x) has a singularity at some point a, then w(x) may not have
that singularity if v(x) tends to 0 as x → a sufficiently fast. It is vital,
of course, that lim v(x) is 0. If it is some non-zero limit L, then in a
x→a
neighbourhood of a, v(x) 6= 0 and so we can recover u(x) from w(x)
as u(x) = w(x)
v(x)
and regularity of w(x) would imply that of u(x) at a, a
contradiction.)
So, we first identify the irregularities of the factor [4x]. It is a function
which is piecewise constant on intervals of length 41 . Its graph is shown
below.

47
y
.

3 [ )

2 [ )

1 [ )

O 1/2 3/4 1 x
1/4

We can now rewrite the function f (x) as




 0 0 < x < 41

 (x − 14 )2 (x − 21 ) 1
4
≤ x < 12
f (x) =  (1)

 2(x − 41 )2 (x − 21 ) 1
2
≤ x < 34
3(x − 41 )2 (x − 21 ) 3

4
≤x<1

It is immediate that in the open interval (0, 1), the function [4x] is
discontinuous only at the three points 14 , 21 and 34 .
Let us now see to what extent these irregularities are cured by the other
factors of f (x), viz. (x − 14 )2 and (x − 12 ). They cure the discontinuities
at 14 and 21 respectively. (For a rigorous proof, [4x] is bounded on
(0, 1). Hence as x → 12 , [4x](x − 21 ) → 0. Similarly, [4x](x − 41 )2 → 0
as x → 14 .) However, discontinuity at 34 remains since the left handed
and right handed limits of f (x) at 43 are 16 2 3
and 16 respectively. Hence
(A) is true.
But the factor (x − 14 )2 is stronger and not only cures discontinuity of
[4x] at x = 41 , but does more. It makes f (x) differentiable at 14 . This is
so because from the first and the second line of the R.H.S. of (1), both
the left handed and the right handed derivatives of f (x) vanish at 14 .
However, f (x) is not differentiable at 21 because from the second line
on the R.H.S. of (1), the left handed derivative of f at 12 , i.e. f−′ ( 21 )
1
equals 16 while from the third line, f+′ ( 21 ) equals 16
2
. (We got these,

48
respectively, by taking the left handed and the right handed limits of
f ′ (x) as x → 21 . Conceptually, they differ from f−′ ( 21 ) and f+′ ( 12 ) But
using Lagrange MVT it can be shown that they are equal.)
So, 12 is a point where f is continuous but not differentiable. At every
other point of (0, 1) f (x) is either differentiable or discontinuous. Hence
(B) is true. As there are only two points, viz. 21 and 34 where f (x) fails
to be differentiable, (C) is false.
Finally, for (D), we note from (1) that f (x) is negative for 14 < x < 12
and positive for x > 12 . So the minimum can occur only on ( 14 , 21 ). On
this interval,
1 1 1
f ′ (x) = 2(x − )(x − ) + (x − )2
4 2 4
1 5
= (x − )(3x − ) (2)
4 4
Hence in ( 14 , 21 ), f ′ vanishes only when x = 12
5
. So the minimum occurs
5 1
at x = 12 . By a direct calculation, its value is − 432 . Hence (D) is false.

A good problem testing the ability to recognise which points of


possible irregularities can be smoothened by multiplying by suitable
factors. Although we have given the reasoning in detail, it does not
take much time to apply it. The last statement (D) is totally out of
place.

Q.7 Let S be the set of all twice differentiable functions f from IR to IR


d2 f
such that 2 > 0 for all x ∈ (−1, 1). For f ∈ X, let Xf be the number
dx
of points x ∈ (−1, 1) for which f (x) = x. Then which of the following
statements is(are) true?

(A) There exists a function f ∈ S such that Xf = 0


(B) For every function f ∈ S, we have Xf ≤ 2
(C) There exists a function f ∈ S such that Xf = 2
(D) There does NOT exist any function f in S such that Xf = 1

49
Answer and Comments: (A, B, C). The condition on f means that it
is strictly concave upwards. Hence the chord always lies above the graph. A
point x for which f (x) = x is often called a fixed point of f . The question
deals with the possible number of fixed points of a strictly concave upward
function.
f (x) = ex is an example of a strictly concave upwards function whose
graph is always above the line y = x as shown in (a). Hence it has no
fixed points. But if you pull it down by 1 unit, i.e. consider the function
f (x) = ex − 1, then it touches the line y = x at the origin as in (b). So there
is exactly one fixed point. If the graph is pulled down still further as in (c),
you get an example of a function with two fixed points.
y y y

y=ex
y=x

x x x

(a) (b) (c)

These three examples show that (A) and (C) are true while (D) is false.
For (B), a geometric argument is that if there are three fixed points, say
a, b, c with a < b < c, then (a, a), (b, b) and (c, c) all lie on the graph of f .
But that contradicts that the chord joining (a, f (a)) and (c, f (c)) lies strictly
above the graph, which in particular implies that f (b) < b. For an analytical
argument, consider a new function g(x) defined by

g(x) = f (x) − x (1)

Fixed points of f are precisely the zeros of g. If f has more than two fixed
points, then g has more than two zeros and hence by Rolle’s theorem, g ′ has
at least two zeros. But g ′′ (x) = f ′′ (x) > 0 everywhere. That means g ′(x) is
strictly increasing and hence cannot have two zeros. That is a contradiction.
So, (B) is true.

A good problem on theoretical calculus, requiring hardly any compu-


tations.

50
SECTION - 3
This section contains SIX questions, the answer to each of which is a non-
negative integer. There are 4 points for a correct answer and 0 in all other
cases.
 
Q.8 For x ∈ IR, let tan−1 (x) ∈ − π2 , π2 . Then the minimum value of the
Z x tan−1 x (t−cos t)
e
function f : IR −→ IR f (x) = dt is
0 1 + t2023
Answer and Comments: 0. x tan−1 x is positive for all x 6= 0. Also
the integrand is positive for all t ≥ 0. Hence the integral is positive for
all x 6= 0. Since it vanishes for x = 0, its minimum value is 0.

An absolutely simple problem where there is no need to evaluate


the integral. Probably the idea is to test the quality of weeding out
pieces of data not relevant to the answer.
Q.9 For x ∈ IR, let y(x) be a solution of the differential equation (x2 −
dy
5) − 2xy = −2x(x2 − 5)2 such that y(2) = 7. Then the maximum
dx
value of the function y(x) is

Answer and Comments: 16. A standard type linear differential


equation of order 1. Rewrite it as
dy 2x
− 2 y = −2x(x2 − 5) (1)
dx x − 5
In Ra neighbourhood of 2, x2 − 5 < 0. So the integrating factor is not
2x
− dx 2 2
e x2 −5 = e− ln(x −5) = x21−5 , but rather, e− ln(5−x ) = 5−x
1
2 . Multiply-

ing by it, the general solution comes to be


Z
y
= 2xdx (2)
5 − x2
or
y = (5 − x2 )(x2 + c) (3)
where c is a constant. The initial condition y(2) = 7 determines c as 3.
Hence the particular solution is
y = y(x) = (5 − x2 )(x2 + 3) (4)

51
To maximise this, we note that it is a biquadratic in x. Putting x2 = u,
we might as well maximise
y(u) = (5 − u)(u + 3) (5)
for u ≥ 0. This can be done by expanding y(u) and taking its derivative.
But a better way is to note that it is a quadratic in u with negative
leading coefficient. So its graph is a downward parabola which meets
the u-axis at 5 and −3. Hence its maximum occurs at the midpoint
u = 1 and the maximum value is 16.

Another straightforward problem. Even if the derivation of the


I.F. through ln(x2 − 5) is mistaken, x21−5 is also a valid I.F. as one can
check by multiplying the d.e. throughout by it. This is not surprising,
because when some function, say g(x), is an I.F. for a d.e. so is −g(x),
and indeed kg(x) for any constant k. As a result, in the present problem
the mistake is masked, and being an MCQ, it is impossible to tell if
the candidate got an advantage of this masking. Already there is Q.1
where the linear differential equation first had to be obtained and then
solved. There was no need to ask the present question on its backdrop.
Q.10 Let X be the set of all five digit numbers formed using 1, 2, 2, 2, 4, 4, 0.
For example, 22240 is in X while 02244 and 44422 are not in X .
Suppose that each element of X has an equal chance of being chosen.
Let p be the conditional probability that an element chosen at random
is a multiple of 20 given that it is a multiple of 5. Then the value of
38p is equal to

Answer and Comments: 31. We want only the conditional prob-


ability of divisibility by 20 given the divisibility by 5. So instead of
finding the size of the entire sample space X, we count only how many
members of it are divisible 5, i.e. how many end with the digit 0, as
this is the only possibility.
A number from X that ends in 0 is of the form x1 x2 x3 x4 0 where
x1 , x2 , x3 , x4 belong to the multiset 1, 2, 2, 2, 4, 4 of six elements. We
classify such numbers into the following types.

(i) (1,2,2,2)-type. (That means a permutation of one 1 and three 2’s.)


There are 4 such members.

52
(ii) (4,2,2,2)-type. Here also there are 4 members.
4!
(iii) (1,4,2,2)-type. There are = 12 such members.
2!
4!
(iv) (2,2,4,4)-type . There are = 6 such members.
2!2!
4!
(v) (1,2,4,4)-type. There are = 12 such members.
2!
Adding, there are in all 4 + 4 + 12 + 6 + 12 = 38 numbers which are
divisible by 5.
We now need to count how many of these are divisible by 20. As the
number is already known to end with a 0, it will be divisible by 20 if
and only if its last but one digit is even. In the present case, this means
that the last but one digit is 2 or 4. As there are many ways this can
happen, we consider complementary counting, i.e. how many numbers
ending with 0 have 1 in the last but one digit, the other restrictions
remaining the same.
When the number ends with 10, there are several possibilities depend-
ing on what the first three digits are.

(a) (2,2,2)-type. Only one number of this type.


(b) (2,2,4)-type. There are 3 numbers of this type.
(c) (2,4,4)-type. There are 3 numbers of this type too.

Thus out of the 38 numbers that are divisible by 5, there are 1+3+3 = 7
that are not divisible by 20. The remaining 31 are divisible by 20.
Hence the conditional probability p that an element chosen at random
is divisible by 20 given that it is a multiple of 5 is 31
38
. Hence 38p = 31.

A good combination of elementary number theory and permuta-


tions. The real test is in classifying the various possibilities for the first
part and then resorting to complementary counting for the second. A
candidate is likely to miss or duplicate some of the possibilities. To
help such candidates, the paper-setters have given an implied hint that
the number of members of X divisible by 5 is 38 (or perhaps 19 which
is a divisor of 38).

53
Q.11 Let A1 , A2 , A3 , . . . , A8 be the vertices of a regular octagon that lie on a
circle of radius 2. Let P be a point on the circle and let P Ai denote the
distance between the points P and Ai for i = 1, 2, . . . , 8. If P varies over
the circle, then the maximum value of the product P A1 · P A2 . . . P A8
is

Answer and Comments: 512. Taking the origin at the centre of


the octagon and the axes suitably, we may suppose that A1 , A2 , . . . , A8
are of the form 2, 2α, 2α2, . . . , 2α7 respectively, where α = e2πi/8 is a
complex 8-th root of unity. Factorising z 8 − 1 we have

z 8 − 1 = (z − 1)(z − α) . . . (z − α7 ) (1)
z
for any complex number z. Replacing z by 2
we get

z 8 − 28 = (z − 2)(z − 2α) . . . (z − 2α7 ) (2)

for every complex number z. If we let z represent the point P then the
absolute value of the R.H.S. is precisely the product P A1 P A2 . . . P A8
which we want to maximise. If we let w = 2z , then (2) becomes

P A1 P A2 . . . P A8 = 28 |w 8 − 1| (3)

Since |z| = 2, |w| = 1. Hence w and therefore w 8 also lies on the unit
circle. Hence the maximum value of |w 8 − 1| is 2. It occurs when w 8 is
diametrically opposite to 1 on the unit circle (i.e. w 8 = −1). Therefore
the maximum value of the L.H.S. is 29 = 512.

Problems based on complex roots of unity are common in JEE. The


present one is a slight modification where we take twice the 8-th roots
of unity.
  

a 3 b 


Q.12 Let R = c 2 d 
  : a, b, c, d ∈ {0, 3, 5, 7, 11, 13, 17, 19} . Then
 

0 5 0 
the number of invertible matrices in R is

Answer and Comments: 3780. It is a favourite game of the JEE


paper-setters to combine elementary linear algebra with elementary

54
counting to come up with problems where the entries of matrices are
restricted to lie in some finite set, say S. Often S of the form {0, 1} or
{0, 1, −1}. It is probably for the first time that we see a set S with 8
elements. The entries of S may seem arbitrary. But if we note that all
except 0 among them are primes, there is an implied hint that prime
factorisation has to be used somewhere.
Let A be a 3 × 3 matrix of the type in the problem. Expansion
w.r.t. the last row gives

|A| = 5(bc − ad) (1)

A will be invertible if and only if ad − bc 6= 0. It is easier to calculate


when this equation holds, i.e. when A is singular. So we resort to
complementary counting. Clearly

|R| = 84 (2)

because each of the entries a, b, c, d can be chosen in 8 ways indepen-


dently of each other.
We now analyse how many solutions the equation

ad = bc (3)

has where a, b, c, d ∈ {0, 3, 5, 7, 11, 13, 17, 19}. If both the sides are
non-zero, then because of unique prime factorisation, the sets {a, d}
and {b, c} must be the same. If a = d, then b and c must also equal a
(=d) because of unique factorisation. That is, a = b = c = d 6= 0. This
can happen in 7 ways. On the other hand, if a 6= d, then the ordered
pair (a, d) can be chosen in 7 × 6 = 42 ways while for each such choice
there are two possibilities, viz. b = a (and correspondingly c = d) or
b = d (with c = a). Hence there are 84 ways in which ad = bc 6= 0
with a 6= d. In all there are 7 + 84 = 91 singular matrices for which
ad = bc 6= 0.
We still have to count those matrices for which ad and bc both vanish.
Vanishing of ad means at least one of a and d is 0. By the principle of
inclusion and exclusion, this happens in 8 + 8 − 1 = 15 ways. Similarly
bc = 0 can happen in 15 ways independently of the values of a and d.
Together, ad = bc = 0 can hold in 15 × 15 = 225 ways.

55
So, in all there are 91 + 225 = 316 ways in which A is singular. As the
total number of matrices is 84 = 4096, the number of non-singular (or
invertible) matrices in R is 4096 − 316 = 3780.

There is a little bit of everything, invertibility of a matrix, its


determinant, unique prime factorisation and the principle of inclusion
and exclusion. Instead of giving 3, 5, 7, 11, 13, 17 and 19, had the peper-
setters given p1 , p2 , p3 , p4 , p5 , p6 , p7 as some distinct primes, that would
have given an unwarranted hint. It is not clear what is gained by mak-
ing S a set with 8 elements. If it had only 0 and some 4 or 5 primes
in it, the reasoning would be the same but the calculations a little less
prone to slips of hand.
Q.13 Let C1 be the circle of radius 1 with center at the origin. Let C2 be the
circle of radius r with center at the point A = (4, 1), where 1 < r < 3.
Two distinct common tangents P Q and ST of C1 and C2 are drawn.
The tangent P Q touches C1 at P and C2 at Q . The tangent ST
touches C1 at S and C2 at T . Mid points of the line segments P Q and
ST are√joined to form a line which meets the x-axis at a point B . If
AB = 5, then the value of r 2 is

Answer and Comments: 2. The data is shown in the figure below


where M and N are the midpoints of the segments P Q and ST respec-
tively. It is tacitly assumed that the two tangents are direct and not
transverse, that is, both the circles lie on the same side of each tangent.
y

Q
M
P A C2
r (4,1)
(1,0)
O B (2,0) x
C1 T
S N

56

If we can get hold of the point B in terms of r, then AB = 5 will give
us an equation in r, solving which we can get the value of r and hence
of r 2 . But to get hold of B, we must first get hold of the line MN. For
that we must get hold of the points M and N and for that we must get
hold of the points of contact P, Q, S, T and for that we must get hold
of the equations of these tangents.
There are formulas for getting the equations of the tangents and
their points of contact, in terms of the equations of the circles. Still,
this is a long drawn process and so we look for a better way to get hold
of the line MN without finding the points M and N first.
Fortunately, there is a radically different way to directly write
down the equation of MN from the equations of C1 and C2 . The
pun is intended because it is based on the concept of what is called
the radical axis of a pair of circles. That concept is again based on
a powerful concept (pun intended again) of the power of a point X
w.r.t. a circle C. If a line through X cuts C in two points, say E and
F , then it can be shown that the product XE.XF is independent of
the line. It is called the power of X w.r.t. C. When the line L is a
tangent to C, the product is the square of the tangent from X to C.
It can be shown that all points which have equal powers w.r.t. two
circles C1 and C2 are collinear and the line on which they lie is called
the radical axis of the two circles C1 and C2 .
If we apply this, then the powers of M w.r.t. C1 and C2 are MP 2
and MQ2 . They are equal since M is the midpoint of P Q. So M lies
on the radical axis of C1 and C2 . Similarly, N does. So the line MN
is nothing but the radical axis of C1 and C2 .
All this would be of little help if there was no easy way to find
the equation of the radical axis of two given circles. But here again,
the situation is pleasant. If we write the equations of the circles in the
standard form as, say,
x2 + y 2 + 2gx + 2f y + c = 0 (1)
and x2 + y 2 + 2g ′ x + 2f ′ y + c′ = 0 (2)
then the radical axis is simply the linear equation that results if we
subtract (2) from (1), i.e. the equation
2(g − g ′)x + 2(f − f ′ )y + c − c′ = 0 (3)

57
In the present problem the equations of the circles are x2 + y 2 = 1 and
x2 + y 2 − 8x − 2y + 17 = r 2 . Hence the radical axis is

8x + 2y − 18 + r 2 = 0 (4)
18−r 2
This cuts the x-axis at B. Putting y = 0 in (4), x = . Since AB
√ 8
is given to be 5 we have

18 − r 2 2
( ) +1=5 (5)
8
2 √
Simplifying, 18−r
8
= 4 = ±2, which gives 18 − r 2 = ±16. This gives
r 2 = 2 or 34. The second possibility has to be discarded since that
would put C1 inside C2 precluding any common tangents.

Once again, it is impossible to tell how many candidates will realise


and discard the possibility that r 2 = 34. (As mentioned in Comment
No. 10 of Chapter 24, in JEE 1993, a question asked for common
tangents to two circles. But because of a printing mistake, one of them
was completely inside the other!)
The problem seems to be tailor made as an application of the
concept of radical axis. Those not familiar with it can begin by finding
the equations of the common tangents by equating their distances from
the centres of the circles with 1 and r respectively. The points of contact
will be the feet of the perpendiculars from the centres to these common
tangents. Once P, Q, S and T are known (all in terms of r, of course),
getting M, N and finally B can be done with elementary formulas.
Once M is obtained, some work can be spared by noting that the line
MN is perpendicular to the line joining the centres of the circles.
Kshitij Mehta has pointed out a shortcut based directly on equating
the powers of B w.r.t. the two circles, once it is known that it lies on
the radical axis, without actually finding the radical axis. In general
if C is a circle of radius r and B a point outside C at a distance d
from the centre of C, then the power of C w.r.t. C is d2 − r 2 . This
follows by considering the chord through B which contains a diameter
of C. To apply √ this, we first note that since B lies on the x-axis and
at a distance 5 from (4, 1), B is either (2, 0) or (6, 0). The second
possibility is discarded as the two circles must lie on opposite sides of

58
their radical axis. Once B is identified as (2, 0), its powers w.r.t. C1
and C2 are 22 − 12 and 5 − r 2 respectively. Equating the two gives
r 2 = 2.

SECTION-4
This section contains TWO paragraphs. Based on each paragraph, there
are TWO questions. The answer to each question is a NUMERICAL VALUE.
For each question, the answer is to be entered by truncating/rounding the
numerical value to TWO decimal places. There are THREE marks for ONLY
the correct numerical value and 0 marks in all other cases.

PARAGRAPH “I”

Consider an obtuse angled triangle ABC in which the difference between


the largest and the smallest angle is π2 and whose sides are in arithmetic
progression. Suppose that the vertices of the triangle lie on a circle of radius
1.

(The following two questions are based on PARAGRAPH “I”.)

Q.14 Let a be the area of the triangle ABC . Then the value of (64a)2 is

Answer and Comments: 1008.00. Let the smallest angle be α.


Then the largest angle is π2 + α and the third angle is π2 − 2α. So
π π
α< − 2α < + α (1)
2 2
By the sine rule, the sides are proportional to the sines of their opposite
angles. So if the sides are in an A.P. so are the sines of the opposite
angles. That gives
π π
sin α + sin( + α) = 2 sin( − 2α) (2)
2 2
which reduces to
sin α + cos α = 2 cos 2α
= 2(cos2 α − sin2 α)
= 2(cos α + sin α)(cos α − sin α) (3)

59
and further to the trigonometric equation (since for an acute angle α,
the L.H.S. is non-zero)
1
cos α − sin α = (4)
2
We can now determine α and then all the sides (because the circumra-
dius is also given). That will enable us to answer any question about
the triangle. But that may be unnecessarily complicated. If we use the
formula 2R2 sin A sin B sin C for the area of a triangle ABC, then since
R is given, we only need the sines of the angles and there may be a
short-cut for it. Indeed, squaring (4) we get 1 − 2 sin α cos α = 41 which
implies
3
sin 2α = (5)
4√
7
and further cos 2α = (6)
4
Using the formula stated above for the area and R = 1, we get
π π
a = 2 sin α sin( + α) sin( − 2α)
2 2
= 2 sin α cos α cos 2α
= sin√2α cos 2α
3 7
= (7)
16
√ !2
3 7
using (5) and (6). Hence (64a)2 = 64 × = 16 × 9 × 7 = 1008.
16

It is not clear what is gained by asking the value of (64a)2 . If the


purpose was merely to have an integer for an answer, it would have
been served equally well by asking the value of (16a)2 which is less
torturous.
Q.15 Then the inradius of the triangle ABC is

Answer and Comments: 0.25. Since the area is already obtained


in the last question the easiest approach is to use the formula
area
r= (8)
s
60
where r is the inradius and s is the semi-perimeter of the triangle i.e.
half the sum of the lengths of the sides. Since the circumradius is 1,
the lengths of the three sides are 2 sin α, 2 sin( π2 + α) and 2 sin( π2 − 2α).
So,

s = sin α + cos α + cos 2α (9)

From (6) we already know cos 2α. As for sin α + cos α, we use (4) to
get
cos 2α
sin α + cos α =
cos α − sin α

7

7
= 41 = (10)
2
2

Hence
√ √ √
7 7 3 7
s= + = (11)
2 4 4
and finally, from (7) and (8),

3 7
16 1
r= √ = (12)
3 7 4
4

Conversion to decimal notation completes the solution.

Most of the work needed was already done for the last question.
The present question is mostly clerical once you use the right formula
(8).

PARAGRAPH “II”

Consider the 6 × 6 square in the figure. Let A1 , A2 , . . . , A49 be the points of


intersections ( dots in the picture ) in some order. We say that Ai and Aj
are friends if they are adjacent along a row or along a column. Assume that
each point Ai has an equal chance of being chosen.

61
. . . . . . .

. . . . . . .

. . . . . . .

. . . . . . .

. . . . . . .

. . . . . . .

. . . . . . .

(The following two questions are based on PARAGRAPH“II”.)

Q.16 Let pi be the probability that a randomly chosen point has i many
friends, i = 0, 1, 2, 3, 4. Let X be a random variable such that for
i = 0, 1, 2, 3, 4, the probability P (X = i) = pi . Then the value of
7E(X) is

Answer and Comments: 24.00. We have to first find the probability


distribution of X which is given to assume 0, 1, 2, 3 and 4 as possible
values. The first two are actually never attained, because no matter
which of the 49 points you choose, it has at least two friends. So
p0 = p1 = 0. The number of friends of a point is 2, 3 or 4 depending
on whether it is a corner point, an edge point (other than a corner) or
an interior point. There are 4 corner points, 20 edge points and the
remaining 25 points are interior points. So,
4
p2 = (1)
49
20
p3 = (2)
49
25
and p4 = (3)
49
Hence
4
X
E(X) = pi i
i=0

62
4 20 25
= ×2+ ×3+ ×4
49 49 49
8 + 60 + 100 168 24
= = = (4)
49 49 7
Hence 7E(X) = 24.

An absolutely simple problem once you understand what is asked.


This testing has some novelty now because this is first year since statis-
tics is back in JEE syllabus after a long gap. In coming years, there
could be generalisations with an n × n grid instead of a 6 × 6 one. Or
the condition for friendship more complicated than simple adjacency.

Q.17 Two distinct points are chosen randomly out of the points A1 , A2 , . . . , A49 .
Let p be the probability that they are friends. Then the value of 7p is

Answer and Comments: 0.50. The  total number of (unordered)


49
pairs from a set of 49 elements is 2 = 49 × 24. We now have to
count how many of these are friendly pairs. That happens if and only
if the two points are the end points of either a horizontal or a vertical
segment of length 1. The two possibilities are mutually exclusive. As
there are 7 horizontal lines and in each there are 6 segments of length
1, in all there are 42 horizontal friendship pairs. Similarly there are
42 vertical friendships. In all there are 84 friendly pairs. Hence the
desired probability p is
84 1
p= = (5)
49 × 24 14
1
Hence 7p = 2
= 0.50.

Another trivial problem.

63
CONCLUDING REMARKS

Overall, the problems in Paper 2 are simpler than those in Paper 1. The
geometry problem (Q.13) is a novel application of the radical axis of two
circles. But it is doubtful if that topic is now covered in schools as there is a
de-emphasis on pure geometry. Q.7 about fixed points of a concave upward
function is also a good question on theoretical calculus. Q.2 is interesting
because problems on infinitistic probability are not asked every year.
In Paper 1, Q.4 is a novel application of the Sandwich Theorem. Q.10 is
easy once you get the key idea. But to arrive at the key idea after systematic
thinking is quite challenging.
There are two problems on differential equations. Both of them are about
linear differential equations. The duplication could have been avoided by
replacing one of them by applications of differential equations.
Number theory enters in a very elementary way in some questions. Statis-
tics has been introduced after a long time. Possibly to avoid controversies,
the paper-setters have asked very elementary questions.

64

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